Update README.md

README.md

@@ -5,6 +5,18 @@ language:
 tags:
 - diffusion
 ---
+SDXL-Dreamshaper
+[SDXL](https://arxiv.org/abs/2307.01952) consists of an [ensemble of experts](https://arxiv.org/abs/2211.01324) pipeline for latent diffusion: 
+In a first step, the base model is used to generate (noisy) latents, 
+which are then further processed with a refinement model (available here: https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0/) specialized for the final denoising steps.
+Note that the base model can be used as a standalone module.
+
+Alternatively, we can use a two-stage pipeline as follows: 
+First, the base model is used to generate latents of the desired output size. 
+In the second step, we use a specialized high-resolution model and apply a technique called SDEdit (https://arxiv.org/abs/2108.01073, also known as "img2img") 
+to the latents generated in the first step, using the same prompt. This technique is slightly slower than the first one, as it requires more function evaluations.
+
+Source code is available at https://github.com/Stability-AI/generative-models .
 ```
 import torch
 from diffusers import StableDiffusionXLPipeline