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+ Basketball is a team sport in which two teams, most commonly of five players each, opposing one another on a rectangular court, compete with the primary objective of shooting a basketball (approximately 9.4 inches (24 cm) in diameter) through the defender's hoop (a basket 18 inches (46 cm) in diameter mounted 10 feet (3.048 m) high to a backboard at each end of the court) while preventing the opposing team from shooting through their own hoop. A field goal is worth two points, unless made from behind the three-point line, when it is worth three. After a foul, timed play stops and the player fouled or designated to shoot a technical foul is given one or more one-point free throws. The team with the most points at the end of the game wins, but if regulation play expires with the score tied, an additional period of play (overtime) is mandated.
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+ Players advance the ball by bouncing it while walking or running (dribbling) or by passing it to a teammate, both of which require considerable skill. On offense, players may use a variety of shots—the lay-up, the jump shot, or a dunk; on defense, they may steal the ball from a dribbler, intercept passes, or block shots; either offense or defense may collect a rebound, that is, a missed shot that bounces from rim or backboard. It is a violation to lift or drag one's pivot foot without dribbling the ball, to carry it, or to hold the ball with both hands then resume dribbling.
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+ The five players on each side fall into five playing positions. The tallest player is usually the center, the second tallest and strongest is the power forward, a slightly shorter but more agile player is the small forward, and the shortest players or the best ball handlers are the shooting guard and the point guard, who implements the coach's game plan by managing the execution of offensive and defensive plays (player positioning). Informally, players may play three-on-three, two-on-two, and one-on-one.
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+ Invented in 1891 by Canadian-American gym teacher James Naismith in Springfield, Massachusetts, United States, basketball has evolved to become one of the world's most popular and widely viewed sports.[1] The National Basketball Association (NBA) is the most significant professional basketball league in the world in terms of popularity, salaries, talent, and level of competition.[2][3] Outside North America, the top clubs from national leagues qualify to continental championships such as the EuroLeague and the Basketball Champions League Americas. The FIBA Basketball World Cup and Men's Olympic Basketball Tournament are the major international events of the sport and attract top national teams from around the world. Each continent hosts regional competitions for national teams, like EuroBasket and FIBA AmeriCup.
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+ The FIBA Women's Basketball World Cup and Women's Olympic Basketball Tournament feature top national teams from continental championships. The main North American league is the WNBA (NCAA Women's Division I Basketball Championship is also popular), whereas strongest European clubs participate in the EuroLeague Women.
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+ In early December 1891, Canadian James Naismith,[4] a physical education professor and instructor at the International Young Men's Christian Association Training School[5] (YMCA) (today, Springfield College) in Springfield, Massachusetts, was trying to keep his gym class active on a rainy day. He sought a vigorous indoor game to keep his students occupied and at proper levels of fitness during the long New England winters. After rejecting other ideas as either too rough or poorly suited to walled-in gymnasiums, he wrote the basic rules and nailed a peach basket onto an elevated track. In contrast with modern basketball nets, this peach basket retained its bottom, and balls had to be retrieved manually after each "basket" or point scored; this proved inefficient, however, so the bottom of the basket was removed, allowing the balls to be poked out with a long dowel each time.
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+ Basketball was originally played with a soccer ball. These round balls from "association football" were made, at the time, with a set of laces to close off the hole needed for inserting the inflatable bladder after the other sewn-together segments of the ball's cover had been flipped outside-in.[6][7] These laces could cause bounce passes and dribbling to be unpredictable.[8] Eventually a lace-free ball construction method was invented, and this change to the game was endorsed by Naismith. (Whereas in American football, the lace construction proved to be advantageous for gripping and remains to this day.) The first balls made specifically for basketball were brown, and it was only in the late 1950s that Tony Hinkle, searching for a ball that would be more visible to players and spectators alike, introduced the orange ball that is now in common use. Dribbling was not part of the original game except for the "bounce pass" to teammates. Passing the ball was the primary means of ball movement. Dribbling was eventually introduced but limited by the asymmetric shape of early balls.[dubious – discuss] Dribbling was common by 1896, with a rule against the double dribble by 1898.[9]
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+ The peach baskets were used until 1906 when they were finally replaced by metal hoops with backboards. A further change was soon made, so the ball merely passed through. Whenever a person got the ball in the basket, his team would gain a point. Whichever team got the most points won the game.[10] The baskets were originally nailed to the mezzanine balcony of the playing court, but this proved impractical when spectators in the balcony began to interfere with shots. The backboard was introduced to prevent this interference; it had the additional effect of allowing rebound shots.[11] Naismith's handwritten diaries, discovered by his granddaughter in early 2006, indicate that he was nervous about the new game he had invented, which incorporated rules from a children's game called duck on a rock, as many had failed before it.
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+ Frank Mahan, one of the players from the original first game, approached Naismith after the Christmas break, in early 1892, asking him what he intended to call his new game. Naismith replied that he hadn't thought of it because he had been focused on just getting the game started. Mahan suggested that it be called "Naismith ball", at which he laughed, saying that a name like that would kill any game. Mahan then said, "Why not call it basketball?" Naismith replied, "We have a basket and a ball, and it seems to me that would be a good name for it."[12][13] The first official game was played in the YMCA gymnasium in Albany, New York, on January 20, 1892, with nine players. The game ended at 1–0; the shot was made from 25 feet (7.6 m), on a court just half the size of a present-day Streetball or National Basketball Association (NBA) court.
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+ At the time, football was being played with 10 to a team (which was increased to 11). When winter weather got too icy to play football, teams were taken indoors, and it was convenient to have them split in half and play basketball with five on each side. By 1897–1898 teams of five became standard.
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+ Basketball's early adherents were dispatched to YMCAs throughout the United States, and it quickly spread through the United States and Canada. By 1895, it was well established at several women's high schools. While YMCA was responsible for initially developing and spreading the game, within a decade it discouraged the new sport, as rough play and rowdy crowds began to detract from YMCA's primary mission. However, other amateur sports clubs, colleges, and professional clubs quickly filled the void. In the years before World War I, the Amateur Athletic Union and the Intercollegiate Athletic Association of the United States (forerunner of the NCAA) vied for control over the rules for the game. The first pro league, the National Basketball League, was formed in 1898 to protect players from exploitation and to promote a less rough game. This league only lasted five years.
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+ James Naismith was instrumental in establishing college basketball. His colleague C.O. Beamis fielded the first college basketball team just a year after the Springfield YMCA game at the suburban Pittsburgh Geneva College.[14] Naismith himself later coached at the University of Kansas for six years, before handing the reins to renowned coach Forrest "Phog" Allen. Naismith's disciple Amos Alonzo Stagg brought basketball to the University of Chicago, while Adolph Rupp, a student of Naismith's at Kansas, enjoyed great success as coach at the University of Kentucky. On February 9, 1895, the first intercollegiate 5-on-5 game was played at Hamline University between Hamline and the School of Agriculture, which was affiliated with the University of Minnesota.[15][16][17] The School of Agriculture won in a 9–3 game.
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+ In 1901, colleges, including the University of Chicago, Columbia University, Cornell University, Dartmouth College, the University of Minnesota, the U.S. Naval Academy, the University of Colorado and Yale University began sponsoring men's games. In 1905, frequent injuries on the football field prompted President Theodore Roosevelt to suggest that colleges form a governing body, resulting in the creation of the Intercollegiate Athletic Association of the United States (IAAUS). In 1910, that body would change its name to the National Collegiate Athletic Association (NCAA). The first Canadian interuniversity basketball game was played at YMCA in Kingston, Ontario on February 6, 1904, when McGill University—Naismith's alma mater—visited Queen's University. McGill won 9–7 in overtime; the score was 7–7 at the end of regulation play, and a ten-minute overtime period settled the outcome. A good turnout of spectators watched the game.[18]
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+ The first men's national championship tournament, the National Association of Intercollegiate Basketball tournament, which still exists as the National Association of Intercollegiate Athletics (NAIA) tournament, was organized in 1937. The first national championship for NCAA teams, the National Invitation Tournament (NIT) in New York, was organized in 1938; the NCAA national tournament would begin one year later. College basketball was rocked by gambling scandals from 1948 to 1951, when dozens of players from top teams were implicated in match fixing and point shaving. Partially spurred by an association with cheating, the NIT lost support to the NCAA tournament.
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+ Before widespread school district consolidation, most American high schools were far smaller than their present-day counterparts. During the first decades of the 20th century, basketball quickly became the ideal interscholastic sport due to its modest equipment and personnel requirements. In the days before widespread television coverage of professional and college sports, the popularity of high school basketball was unrivaled in many parts of America. Perhaps the most legendary of high school teams was Indiana's Franklin Wonder Five, which took the nation by storm during the 1920s, dominating Indiana basketball and earning national recognition.
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+ Today virtually every high school in the United States fields a basketball team in varsity competition.[19] Basketball's popularity remains high, both in rural areas where they carry the identification of the entire community, as well as at some larger schools known for their basketball teams where many players go on to participate at higher levels of competition after graduation. In the 2016–17 season, 980,673 boys and girls represented their schools in interscholastic basketball competition, according to the National Federation of State High School Associations.[20] The states of Illinois, Indiana and Kentucky are particularly well known for their residents' devotion to high school basketball, commonly called Hoosier Hysteria in Indiana; the critically acclaimed film Hoosiers shows high school basketball's depth of meaning to these communities.
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+ There is currently no tournament to determine a national high school champion. The most serious effort was the National Interscholastic Basketball Tournament at the University of Chicago from 1917 to 1930. The event was organized by Amos Alonzo Stagg and sent invitations to state champion teams. The tournament started out as a mostly Midwest affair but grew. In 1929 it had 29 state champions. Faced with opposition from the National Federation of State High School Associations and North Central Association of Colleges and Schools that bore a threat of the schools losing their accreditation the last tournament was in 1930. The organizations said they were concerned that the tournament was being used to recruit professional players from the prep ranks.[21] The tournament did not invite minority schools or private/parochial schools.
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+ The National Catholic Interscholastic Basketball Tournament ran from 1924 to 1941 at Loyola University.[22] The National Catholic Invitational Basketball Tournament from 1954 to 1978 played at a series of venues, including Catholic University, Georgetown and George Mason.[23] The National Interscholastic Basketball Tournament for Black High Schools was held from 1929 to 1942 at Hampton Institute.[24] The National Invitational Interscholastic Basketball Tournament was held from 1941 to 1967 starting out at Tuskegee Institute. Following a pause during World War II it resumed at Tennessee State College in Nashville. The basis for the champion dwindled after 1954 when Brown v. Board of Education began an integration of schools. The last tournaments were held at Alabama State College from 1964 to 1967.[25]
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+ Teams abounded throughout the 1920s. There were hundreds of men's professional basketball teams in towns and cities all over the United States, and little organization of the professional game. Players jumped from team to team and teams played in armories and smoky dance halls. Leagues came and went. Barnstorming squads such as the Original Celtics and two all-African American teams, the New York Renaissance Five ("Rens") and the (still existing) Harlem Globetrotters played up to two hundred games a year on their national tours.
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+ In 1946, the Basketball Association of America (BAA) was formed. The first game was played in Toronto, Ontario, Canada between the Toronto Huskies and New York Knickerbockers on November 1, 1946. Three seasons later, in 1949, the BAA merged with the National Basketball League (NBL) to form the National Basketball Association (NBA). By the 1950s, basketball had become a major college sport, thus paving the way for a growth of interest in professional basketball. In 1959, a basketball hall of fame was founded in Springfield, Massachusetts, site of the first game. Its rosters include the names of great players, coaches, referees and people who have contributed significantly to the development of the game. The hall of fame has people who have accomplished many goals in their career in basketball. An upstart organization, the American Basketball Association, emerged in 1967 and briefly threatened the NBA's dominance until the ABA-NBA merger in 1976. Today the NBA is the top professional basketball league in the world in terms of popularity, salaries, talent, and level of competition.
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+ The NBA has featured many famous players, including George Mikan, the first dominating "big man"; ball-handling wizard Bob Cousy and defensive genius Bill Russell of the Boston Celtics; charismatic center Wilt Chamberlain, who originally played for the barnstorming Harlem Globetrotters; all-around stars Oscar Robertson and Jerry West; more recent big men Kareem Abdul-Jabbar, Shaquille O'Neal, Hakeem Olajuwon and Karl Malone; playmakers John Stockton, Isiah Thomas and Steve Nash; crowd-pleasing forwards Julius Erving and Charles Barkley; European stars Dirk Nowitzki, Pau Gasol and Tony Parker; more recent superstars LeBron James, Allen Iverson and Kobe Bryant; and the three players who many credit with ushering the professional game to its highest level of popularity during the 1980s and 1990s: Larry Bird, Earvin "Magic" Johnson, and Michael Jordan.
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+ In 2001, the NBA formed a developmental league, the National Basketball Development League (later known as the NBA D-League and then the NBA G League after a branding deal with Gatorade). As of the 2018–19 season, the G League has 27 teams.
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+ FIBA (International Basketball Federation) was formed in 1932 by eight founding nations: Argentina, Czechoslovakia, Greece, Italy, Latvia, Portugal, Romania and Switzerland. At this time, the organization only oversaw amateur players. Its acronym, derived from the French Fédération Internationale de Basket-ball Amateur, was thus "FIBA". Men's basketball was first included at the Berlin 1936 Summer Olympics, although a demonstration tournament was held in 1904. The United States defeated Canada in the first final, played outdoors. This competition has usually been dominated by the United States, whose team has won all but three titles. The first of these came in a controversial final game in Munich in 1972 against the Soviet Union, in which the ending of the game was replayed three times until the Soviet Union finally came out on top.[26] In 1950 the first FIBA World Championship for men, now known as the FIBA Basketball World Cup, was held in Argentina. Three years later, the first FIBA World Championship for women, now known as the FIBA Women's Basketball World Cup, was held in Chile. Women's basketball was added to the Olympics in 1976, which were held in Montreal, Quebec, Canada with teams such as the Soviet Union, Brazil and Australia rivaling the American squads.
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+ In 1989, FIBA allowed professional NBA players to participate in the Olympics for the first time. Prior to the 1992 Summer Olympics, only European and South American teams were allowed to field professionals in the Olympics. The United States' dominance continued with the introduction of the original Dream Team. In the 2004 Athens Olympics, the United States suffered its first Olympic loss while using professional players, falling to Puerto Rico (in a 19-point loss) and Lithuania in group games, and being eliminated in the semifinals by Argentina. It eventually won the bronze medal defeating Lithuania, finishing behind Argentina and Italy. The Redeem Team, won gold at the 2008 Olympics, and the B-Team, won gold at the 2010 FIBA World Championship in Turkey despite featuring no players from the 2008 squad. The United States continued its dominance as they won gold at the 2012 Olympics, 2014 FIBA World Cup and the 2016 Olympics.
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+ Worldwide, basketball tournaments are held for boys and girls of all age levels. The global popularity of the sport is reflected in the nationalities represented in the NBA. Players from all six inhabited continents currently play in the NBA. Top international players began coming into the NBA in the mid-1990s, including Croatians Dražen Petrović and Toni Kukoč, Serbian Vlade Divac, Lithuanians Arvydas Sabonis and Šarūnas Marčiulionis, Dutchman Rik Smits and German Detlef Schrempf.
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+ In the Philippines, the Philippine Basketball Association's first game was played on April 9, 1975 at the Araneta Coliseum in Cubao, Quezon City. Philippines. It was founded as a "rebellion" of several teams from the now-defunct Manila Industrial and Commercial Athletic Association, which was tightly controlled by the Basketball Association of the Philippines (now defunct), the then-FIBA recognized national association. Nine teams from the MICAA participated in the league's first season that opened on April 9, 1975. The NBL is Australia's pre-eminent men's professional basketball league. The league commenced in 1979, playing a winter season (April–September) and did so until the completion of the 20th season in 1998. The 1998–99 season, which commenced only months later, was the first season after the shift to the current summer season format (October–April). This shift was an attempt to avoid competing directly against Australia's various football codes. It features 8 teams from around Australia and one in New Zealand. A few players including Luc Longley, Andrew Gaze, Shane Heal, Chris Anstey and Andrew Bogut made it big internationally, becoming poster figures for the sport in Australia. The Women's National Basketball League began in 1981.
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+ Women's basketball began in 1892 at Smith College when Senda Berenson, a physical education teacher, modified Naismith's rules for women. Shortly after she was hired at Smith, she went to Naismith to learn more about the game.[27] Fascinated by the new sport and the values it could teach, she organized the first women's collegiate basketball game on March 21, 1893, when her Smith freshmen and sophomores played against one another.[28] However, the first women's interinstitutional game was played in 1892 between the University of California and Miss Head's School.[29] Berenson's rules were first published in 1899, and two years later she became the editor of A. G. Spalding's first Women's Basketball Guide.[28] Berenson's freshmen played the sophomore class in the first women's intercollegiate basketball game at Smith College, March 21, 1893.[30] The same year, Mount Holyoke and Sophie Newcomb College (coached by Clara Gregory Baer) women began playing basketball. By 1895, the game had spread to colleges across the country, including Wellesley, Vassar, and Bryn Mawr. The first intercollegiate women's game was on April 4, 1896. Stanford women played Berkeley, 9-on-9, ending in a 2–1 Stanford victory.
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+ Women's basketball development was more structured than that for men in the early years. In 1905, the Executive Committee on Basket Ball Rules (National Women's Basketball Committee) was created by the American Physical Education Association.[31] These rules called for six to nine players per team and 11 officials. The International Women's Sports Federation (1924) included a women's basketball competition. 37 women's high school varsity basketball or state tournaments were held by 1925. And in 1926, the Amateur Athletic Union backed the first national women's basketball championship, complete with men's rules.[31] The Edmonton Grads, a touring Canadian women's team based in Edmonton, Alberta, operated between 1915 and 1940. The Grads toured all over North America, and were exceptionally successful. They posted a record of 522 wins and only 20 losses over that span, as they met any team that wanted to challenge them, funding their tours from gate receipts.[32] The Grads also shone on several exhibition trips to Europe, and won four consecutive exhibition Olympics tournaments, in 1924, 1928, 1932, and 1936; however, women's basketball was not an official Olympic sport until 1976. The Grads' players were unpaid, and had to remain single. The Grads' style focused on team play, without overly emphasizing skills of individual players. The first women's AAU All-America team was chosen in 1929.[31] Women's industrial leagues sprang up throughout the United States, producing famous athletes, including Babe Didrikson of the Golden Cyclones, and the All American Red Heads Team, which competed against men's teams, using men's rules. By 1938, the women's national championship changed from a three-court game to two-court game with six players per team.[31]
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+ The NBA-backed Women's National Basketball Association (WNBA) began in 1997. Though it had shaky attendance figures, several marquee players (Lisa Leslie, Diana Taurasi, and Candace Parker among others) have helped the league's popularity and level of competition. Other professional women's basketball leagues in the United States, such as the American Basketball League (1996–98), have folded in part because of the popularity of the WNBA. The WNBA has been looked at by many as a niche league. However, the league has recently taken steps forward. In June 2007, the WNBA signed a contract extension with ESPN. The new television deal ran from 2009 to 2016. Along with this deal, came the first ever rights fees to be paid to a women's professional sports league. Over the eight years of the contract, "millions and millions of dollars" were "dispersed to the league's teams." In a March 12, 2009 article, NBA commissioner David Stern said that in the bad economy, "the NBA is far less profitable than the WNBA. We're losing a lot of money among a large number of teams. We're budgeting the WNBA to break even this year."[33]
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+ Measurements and time limits discussed in this section often vary among tournaments and organizations; international and NBA rules are used in this section.
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+ The object of the game is to outscore one's opponents by throwing the ball through the opponents' basket from above while preventing the opponents from doing so on their own. An attempt to score in this way is called a shot. A successful shot is worth two points, or three points if it is taken from beyond the three-point arc 6.75 metres (22 ft 2 in) from the basket in international games[citation needed] and 23 feet 9 inches (7.24 m) in NBA games.[34] A one-point shot can be earned when shooting from the foul line after a foul is made. After a team has scored from a field goal or free throw, play is resumed with a throw-in awarded to the non-scoring team taken from a point beyond the endline of the court where the points(s) were scored.[35]
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+ Games are played in four quarters of 10 (FIBA)[36] or 12 minutes (NBA).[37] College men's games use two 20-minute halves,[38] college women's games use 10-minute quarters,[39] and most United States high school varsity games use 8-minute quarters; however, this varies from state to state.[40][41] 15 minutes are allowed for a half-time break under FIBA, NBA, and NCAA rules[38][42][43] and 10 minutes in United States high schools.[40] Overtime periods are five minutes in length[38][44][45] except for high school, which is four minutes in length.[40] Teams exchange baskets for the second half. The time allowed is actual playing time; the clock is stopped while the play is not active. Therefore, games generally take much longer to complete than the allotted game time, typically about two hours.
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+ Five players from each team may be on the court at one time.[46][47][48][49] Substitutions are unlimited but can only be done when play is stopped. Teams also have a coach, who oversees the development and strategies of the team, and other team personnel such as assistant coaches, managers, statisticians, doctors and trainers.
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+ For both men's and women's teams, a standard uniform consists of a pair of shorts and a jersey with a clearly visible number, unique within the team, printed on both the front and back. Players wear high-top sneakers that provide extra ankle support. Typically, team names, players' names and, outside of North America, sponsors are printed on the uniforms.
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+ A limited number of time-outs, clock stoppages requested by a coach (or sometimes mandated in the NBA) for a short meeting with the players, are allowed. They generally last no longer than one minute (100 seconds in the NBA) unless, for televised games, a commercial break is needed.
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+ The game is controlled by the officials consisting of the referee (referred to as crew chief in the NBA), one or two umpires (referred to as referees in the NBA) and the table officials. For college, the NBA, and many high schools, there are a total of three referees on the court. The table officials are responsible for keeping track of each team's scoring, timekeeping, individual and team fouls, player substitutions, team possession arrow, and the shot clock.
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+ The only essential equipment in a basketball game is the ball and the court: a flat, rectangular surface with baskets at opposite ends. Competitive levels require the use of more equipment such as clocks, score sheets, scoreboard(s), alternating possession arrows, and whistle-operated stop-clock systems.
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+ A regulation basketball court in international games is 91.9 feet (28.0 meters) long and 49.2 feet (15 meters) wide. In the NBA and NCAA the court is 94 by 50 feet (29 by 15 meters).[34] Most courts have wood flooring, usually constructed from maple planks running in the same direction as the longer court dimension.[50][51] The name and logo of the home team is usually painted on or around the center circle.
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+ The basket is a steel rim 18 inches (46 cm) diameter with an attached net affixed to a backboard that measures 6 by 3.5 feet (1.8 by 1.1 meters) and one basket is at each end of the court. The white outlined box on the backboard is 18 inches (46 cm) high and 2 feet (61 cm) wide. At almost all levels of competition, the top of the rim is exactly 10 feet (3.05 meters) above the court and 4 feet (1.22 meters) inside the baseline. While variation is possible in the dimensions of the court and backboard, it is considered important for the basket to be of the correct height – a rim that is off by just a few inches can have an adverse effect on shooting.
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+ The size of the basketball is also regulated. For men, the official ball is 29.5 inches (75 cm) in circumference (size 7, or a "295 ball") and weighs 22 oz (623.69 grams). If women are playing, the official basketball size is 28.5 inches (72 cm) in circumference (size 6, or a "285 ball") with a weight of 20 oz (567 grams). In 3x3, a formalized version of the halfcourt 3-on-3 game, a dedicated ball with the circumference of a size 6 ball but the weight of a size 7 ball is used in all competitions (men's, women's, and mixed teams).[52]
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+ The ball may be advanced toward the basket by being shot, passed between players, thrown, tapped, rolled or dribbled (bouncing the ball while running).
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+ The ball must stay within the court; the last team to touch the ball before it travels out of bounds forfeits possession. The ball is out of bounds if it touches a boundary line, or touches any player or object that is out of bounds.
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+
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+ There are limits placed on the steps a player may take without dribbling, which commonly results in an infraction known as traveling. Nor may a player stop his dribble and then resume dribbling. A dribble that touches both hands is considered stopping the dribble, giving this infraction the name double dribble. Within a dribble, the player cannot carry the ball by placing his hand on the bottom of the ball; doing so is known as carrying the ball. A team, once having established ball control in the front half of their court, may not return the ball to the backcourt and be the first to touch it. A violation of these rules results in loss of possession.
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+ The ball may not be kicked, nor be struck with the fist. For the offense, a violation of these rules results in loss of possession; for the defense, most leagues reset the shot clock and the offensive team is given possession of the ball out of bounds.
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+
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+ There are limits imposed on the time taken before progressing the ball past halfway (8 seconds in FIBA and the NBA; 10 seconds in NCAA and high school for both sexes), before attempting a shot (24 seconds in FIBA, the NBA, and U Sports (Canadian universities) play for both sexes, and 30 seconds in NCAA play for both sexes), holding the ball while closely guarded (5 seconds), and remaining in the restricted area known as the free-throw lane, (or the "key") (3 seconds). These rules are designed to promote more offense.
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+
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+ Basket interference, or goaltending is a violation charged when a player illegally interferes with a shot. This violation is incurred when a player touches the ball on its downward trajectory to the basket, unless it is obvious that the ball has no chance of entering the basket, if a player touches the ball while it is in the rim, or in the area extended upwards from the basket, or if a player reaches through the basket to interfere with the shot. When a defensive player is charged with goaltending, the basket is awarded. If an offensive player commits the infraction, the basket is cancelled. In either case possession of the ball is turned over to the defensive team.
96
+
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+ An attempt to unfairly disadvantage an opponent through certain types of physical contact is illegal and is called a personal foul. These are most commonly committed by defensive players; however, they can be committed by offensive players as well. Players who are fouled either receive the ball to pass inbounds again, or receive one or more free throws if they are fouled in the act of shooting, depending on whether the shot was successful. One point is awarded for making a free throw, which is attempted from a line 15 feet (4.6 m) from the basket.
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+
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+ The referee is responsible for judging whether contact is illegal, sometimes resulting in controversy. The calling of fouls can vary between games, leagues and referees.
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+
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+ There is a second category of fouls called technical fouls, which may be charged for various rules violations including failure to properly record a player in the scorebook, or for unsportsmanlike conduct. These infractions result in one or two free throws, which may be taken by any of the five players on the court at the time. Repeated incidents can result in disqualification. A blatant foul involving physical contact that is either excessive or unnecessary is called an intentional foul (flagrant foul in the NBA). In FIBA and NCAA women's basketball, a foul resulting in ejection is called a disqualifying foul, while in leagues other than the NBA, such a foul is referred to as flagrant.
102
+
103
+ If a team exceeds a certain limit of team fouls in a given period (quarter or half) – four for NBA, NCAA women's, and international games – the opposing team is awarded one or two free throws on all subsequent non-shooting fouls for that period, the number depending on the league. In the US college men's game and high school games for both sexes, if a team reaches 7 fouls in a half, the opposing team is awarded one free throw, along with a second shot if the first is made. This is called shooting "one-and-one". If a team exceeds 10 fouls in the half, the opposing team is awarded two free throws on all subsequent fouls for the half.
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+
105
+ When a team shoots foul shots, the opponents may not interfere with the shooter, nor may they try to regain possession until the last or potentially last free throw is in the air.
106
+
107
+ After a team has committed a specified number of fouls, the other team is said to be "in the bonus". On scoreboards, this is usually signified with an indicator light reading "Bonus" or "Penalty" with an illuminated directional arrow or dot indicating that team is to receive free throws when fouled by the opposing team. (Some scoreboards also indicate the number of fouls committed.)
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+ If a team misses the first shot of a two-shot situation, the opposing team must wait for the completion of the second shot before attempting to reclaim possession of the ball and continuing play.
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+ If a player is fouled while attempting a shot and the shot is unsuccessful, the player is awarded a number of free throws equal to the value of the attempted shot. A player fouled while attempting a regular two-point shot thus receives two shots, and a player fouled while attempting a three-point shot receives three shots.
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+ If a player is fouled while attempting a shot and the shot is successful, typically the player will be awarded one additional free throw for one point. In combination with a regular shot, this is called a "three-point play" or "four-point play" (or more colloquially, an "and one") because of the basket made at the time of the foul (2 or 3 points) and the additional free throw (1 point).
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+ Although the rules do not specify any positions whatsoever, they have evolved as part of basketball. During the early years of basketball's evolution, two guards, two forwards, and one center were used. In more recent times specific positions evolved, but the current trend, advocated by many top coaches including Mike Krzyzewski is towards positionless basketball, where big guys are free to shoot from outside and dribble if their skill allows it.[53] Popular descriptions of positions include:
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+ Point guard (often called the "1") : usually the fastest player on the team, organizes the team's offense by controlling the ball and making sure that it gets to the right player at the right time.
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+ Shooting guard (the "2") : creates a high volume of shots on offense, mainly long-ranged; and guards the opponent's best perimeter player on defense.
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+ Small forward (the "3") : often primarily responsible for scoring points via cuts to the basket and dribble penetration; on defense seeks rebounds and steals, but sometimes plays more actively.
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+ Power forward (the "4"): plays offensively often with their back to the basket; on defense, plays under the basket (in a zone defense) or against the opposing power forward (in man-to-man defense).
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+ Center (the "5"): uses height and size to score (on offense), to protect the basket closely (on defense), or to rebound.
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+ The above descriptions are flexible. For most teams today, the shooting guard and small forward have very similar responsibilities and are often called the wings, as do the power forward and center, who are often called post players. While most teams describe two players as guards, two as forwards, and one as a center, on some occasions teams choose to call them by different designations.
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+
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+ There are two main defensive strategies: zone defense and man-to-man defense. In a zone defense, each player is assigned to guard a specific area of the court. Zone defenses often allow the defense to double team the ball, a manoeuver known as a trap. In a man-to-man defense, each defensive player guards a specific opponent.
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+ Offensive plays are more varied, normally involving planned passes and movement by players without the ball. A quick movement by an offensive player without the ball to gain an advantageous position is known as a cut. A legal attempt by an offensive player to stop an opponent from guarding a teammate, by standing in the defender's way such that the teammate cuts next to him, is a screen or pick. The two plays are combined in the pick and roll, in which a player sets a pick and then "rolls" away from the pick towards the basket. Screens and cuts are very important in offensive plays; these allow the quick passes and teamwork, which can lead to a successful basket. Teams almost always have several offensive plays planned to ensure their movement is not predictable. On court, the point guard is usually responsible for indicating which play will occur.
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+ Shooting is the act of attempting to score points by throwing the ball through the basket, methods varying with players and situations.
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+ Typically, a player faces the basket with both feet facing the basket. A player will rest the ball on the fingertips of the dominant hand (the shooting arm) slightly above the head, with the other hand supporting the side of the ball. The ball is usually shot by jumping (though not always) and extending the shooting arm. The shooting arm, fully extended with the wrist fully bent, is held stationary for a moment following the release of the ball, known as a follow-through. Players often try to put a steady backspin on the ball to absorb its impact with the rim. The ideal trajectory of the shot is somewhat controversial, but generally a proper arc is recommended. Players may shoot directly into the basket or may use the backboard to redirect the ball into the basket.
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+ The two most common shots that use the above described setup are the set shot and the jump shot. Both are preceeded by a crouching action which preloads the muscles and increases the power of the shot. In a set shot the shooter straightens up and throws from a standing position with neither foot leaving the floor; this is typically used for free throws. For a jump shot, the throw is taken in mid-air with the ball being released near the top of the jump. This provides much greater power and range, and it also allows the player to elevate over the defender. Failure to release the ball before the feet return to the floor is considered a traveling violation.
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+ Another common shot is called the lay-up. This shot requires the player to be in motion toward the basket, and to "lay" the ball "up" and into the basket, typically off the backboard (the backboard-free, underhand version is called a finger roll). The most crowd-pleasing and typically highest-percentage accuracy shot is the slam dunk, in which the player jumps very high and throws the ball downward, through the basket while touching it.
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+ Another shot that is becoming common[citation needed] is the "circus shot". The circus shot is a low-percentage shot that is flipped, heaved, scooped, or flung toward the hoop while the shooter is off-balance, airborne, falling down, and/or facing away from the basket. A back-shot is a shot taken when the player is facing away from the basket, and may be shot with the dominant hand, or both; but there is a very low chance that the shot will be successful.
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+
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+ A shot that misses both the rim and the backboard completely is referred to as an air ball. A particularly bad shot, or one that only hits the backboard, is jocularly called a brick. The hang time is the length of time a player stays in the air after jumping, either to make a slam dunk, lay-up or jump shot.
144
+
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+ The objective of rebounding is to successfully gain possession of the basketball after a missed field goal or free throw, as it rebounds from the hoop or backboard. This plays a major role in the game, as most possessions end when a team misses a shot. There are two categories of rebounds: offensive rebounds, in which the ball is recovered by the offensive side and does not change possession, and defensive rebounds, in which the defending team gains possession of the loose ball. The majority of rebounds are defensive, as the team on defense tends to be in better position to recover missed shots.
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+ A pass is a method of moving the ball between players. Most passes are accompanied by a step forward to increase power and are followed through with the hands to ensure accuracy.
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+ A staple pass is the chest pass. The ball is passed directly from the passer's chest to the receiver's chest. A proper chest pass involves an outward snap of the thumbs to add velocity and leaves the defence little time to react.
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+ Another type of pass is the bounce pass. Here, the passer bounces the ball crisply about two-thirds of the way from his own chest to the receiver. The ball strikes the court and bounces up toward the receiver. The bounce pass takes longer to complete than the chest pass, but it is also harder for the opposing team to intercept (kicking the ball deliberately is a violation). Thus, players often use the bounce pass in crowded moments, or to pass around a defender.
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+ The overhead pass is used to pass the ball over a defender. The ball is released while over the passer's head.
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+ The outlet pass occurs after a team gets a defensive rebound. The next pass after the rebound is the outlet pass.
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+ The crucial aspect of any good pass is it being difficult to intercept. Good passers can pass the ball with great accuracy and they know exactly where each of their other teammates prefers to receive the ball. A special way of doing this is passing the ball without looking at the receiving teammate. This is called a no-look pass.
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+ Another advanced style of passing is the behind-the-back pass, which, as the description implies, involves throwing the ball behind the passer's back to a teammate. Although some players can perform such a pass effectively, many coaches discourage no-look or behind-the-back passes, believing them to be difficult to control and more likely to result in turnovers or violations.
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+ Dribbling is the act of bouncing the ball continuously with one hand and is a requirement for a player to take steps with the ball. To dribble, a player pushes the ball down towards the ground with the fingertips rather than patting it; this ensures greater control.
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+
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+ When dribbling past an opponent, the dribbler should dribble with the hand farthest from the opponent, making it more difficult for the defensive player to get to the ball. It is therefore important for a player to be able to dribble competently with both hands.
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+ Good dribblers (or "ball handlers") tend to bounce the ball low to the ground, reducing the distance of travel of the ball from the floor to the hand, making it more difficult for the defender to "steal" the ball. Good ball handlers frequently dribble behind their backs, between their legs, and switch directions suddenly, making a less predictable dribbling pattern that is more difficult to defend against. This is called a crossover, which is the most effective way to move past defenders while dribbling.
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+ A skilled player can dribble without watching the ball, using the dribbling motion or peripheral vision to keep track of the ball's location. By not having to focus on the ball, a player can look for teammates or scoring opportunities, as well as avoid the danger of having someone steal the ball away from him/her.
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+ A block is performed when, after a shot is attempted, a defender succeeds in altering the shot by touching the ball. In almost all variants of play, it is illegal to touch the ball after it is in the downward path of its arc; this is known as goaltending. It is also illegal under NBA and Men's NCAA basketball to block a shot after it has touched the backboard, or when any part of the ball is directly above the rim. Under international rules it is illegal to block a shot that is in the downward path of its arc or one that has touched the backboard until the ball has hit the rim. After the ball hits the rim, it is again legal to touch it even though it is no longer considered as a block performed.
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+ To block a shot, a player has to be able to reach a point higher than where the shot is released. Thus, height can be an advantage in blocking. Players who are taller and playing the power forward or center positions generally record more blocks than players who are shorter and playing the guard positions. However, with good timing and a sufficiently high vertical leap, even shorter players can be effective shot blockers.
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+ At the professional level, most male players are above 6 feet 3 inches (1.91 m) and most women above 5 feet 7 inches (1.70 m). Guards, for whom physical coordination and ball-handling skills are crucial, tend to be the smallest players. Almost all forwards in the top men's pro leagues are 6 feet 6 inches (1.98 m) or taller. Most centers are over 6 feet 10 inches (2.08 m) tall. According to a survey given to all NBA teams,[when?] the average height of all NBA players is just under 6 feet 7 inches (2.01 m), with the average weight being close to 222 pounds (101 kg). The tallest players ever in the NBA were Manute Bol and Gheorghe Mureșan, who were both 7 feet 7 inches (2.31 m) tall. At 7 feet 2 inches (2.18 m), Margo Dydek was the tallest player in the history of the WNBA.
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+ The shortest player ever to play in the NBA is Muggsy Bogues at 5 feet 3 inches (1.60 m).[54] Other short players have thrived at the pro level. Anthony "Spud" Webb was just 5 feet 7 inches (1.70 m) tall, but had a 42-inch (1.1 m) vertical leap, giving him significant height when jumping. While shorter players are often at a disadvantage in certain aspects of the game, their ability to navigate quickly through crowded areas of the court and steal the ball by reaching low are strengths.
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+ Players regularly inflate their height. Many prospects exaggerate their height while in high school or college to make themselves more appealing to coaches and scouts, who prefer taller players. Charles Barkley stated; "I've been measured at 6-5, 6-4 ​3⁄4. But I started in college at 6-6." Sam Smith, a former writer from The Chicago Tribune, said: "We sort of know the heights, because after camp, the sheet comes out. But you use that height, and the player gets mad. And then you hear from his agent. Or you file your story with the right height, and the copy desk changes it because they have the 'official' N.B.A. media guide, which is wrong. So you sort of go along with the joke."[55] In the NBA, there is no standard on whether a player's listed height uses their measurement with shoes on or without. The NBA Draft Combine, which most players attend before the draft, provides both measurements. Thereafter, a player's team is solely responsible for their listed height, which can vary depending on the process selected.[56][57]
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+
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+ Notable players who overstated their height include:
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+ On rare occasions, some players will understate their actual heights, not to be repositioned. One example is Kevin Durant, whose listed height is 6 feet 9 inches (2.06 m), while his actual height is 7 feet 0 inches (2.13 m). Durant's reasoning was, "Really, that's the prototypical size for a small forward. Anything taller than that, and they'll start saying, 'Ah, he's a power forward."[62]
182
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+ Variations of basketball are activities based on the game of basketball, using common basketball skills and equipment (primarily the ball and basket). Some variations are only superficial rules changes, while others are distinct games with varying degrees of basketball influences. Other variations include children's games, contests or activities meant to help players reinforce skills.
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+
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+ There are principal basketball sports with variations on basketball including Wheelchair basketball, Water basketball, Beach basketball, Slamball, Streetball and Unicycle basketball. An earlier version of basketball, played primarily by women and girls, was Six-on-six basketball. Horseball is a game played on horseback where a ball is handled and points are scored by shooting it through a high net (approximately 1.5m×1.5m). The sport is like a combination of polo, rugby, and basketball. There is even a form played on donkeys known as Donkey basketball, but that version has come under attack from animal rights groups.
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+ Perhaps the single most common variation of basketball is the half-court game, played in informal settings without referees or strict rules. Only one basket is used, and the ball must be "taken back" or "cleared" – passed or dribbled outside the three-point line each time possession of the ball changes from one team to the other. Half-court games require less cardiovascular stamina, since players need not run back and forth a full court. Half-court raises the number of players that can use a court or, conversely, can be played if there is an insufficient number to form full 5-on-5 teams.
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+
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+ Half-court basketball is usually played 1-on-1, 2-on-2 or 3-on-3. The latter variation is gradually gaining official recognition as 3x3, originally known as FIBA 33. It was first tested at the 2007 Asian Indoor Games in Macau and the first official tournaments were held at the 2009 Asian Youth Games and the 2010 Youth Olympics, both in Singapore. The first FIBA 3x3 Youth World Championships[63] were held in Rimini, Italy in 2011, with the first FIBA 3x3 World Championships for senior teams following a year later in Athens. The sport is highly tipped to become an Olympic sport as early as 2016.[64] In the summer of 2017, the BIG3 basketball league, a professional 3x3 half court basketball league that features former NBA players, began. The BIG3 features several rule variants including a four-point field goal.[65]
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+ There are also other basketball sports, such as:
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+
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+ Spin-offs from basketball that are now separate sports include:
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+
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+ Basketball has been adopted by various social groups, which have established their own environments and sometimes their own rules. Such socialized forms of basketball include the following.
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+ Basketball is played widely casually in schools and colleges where fun, entertainment and camaraderie rule rather than winning a game.
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+ Disabled basketball is played by various disabled groups, such as the deaf and physically crippled people.
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+ Show basketball is performed by entertainment basketball show teams, the prime example being the Harlem Globetrotters. There are even specialized entertainment teams, such as teams of celebrities, people with short heights and others.
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+ Fantasy basketball was popularized during the 1990s after the advent of the Internet. Those who play this game are sometimes referred to as General Managers, who draft actual NBA players and compute their basketball statistics. The game was popularized by ESPN Fantasy Sports, NBA.com, and Yahoo! Fantasy Sports. Other sports websites provided the same format keeping the game interesting with participants actually owning specific players.
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1
+
2
+
3
+ Slavery has historically been widespread in Africa, and still continues today in some African countries.
4
+
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+ Systems of servitude and slavery were common in parts of Africa in ancient times, as they were in much of the rest of the ancient world. When the Arab slave trade (which started in the 7th century) and Atlantic slave trade (which started in the 16th century) began, many of the pre-existing local African slave systems began supplying captives for slave markets outside Africa.[1]
6
+
7
+ Slavery in historical Africa was practiced in many different forms: Debt slavery, enslavement of war captives, military slavery, slavery for prostitution, and criminal slavery were all practiced in various parts of Africa.[2] Slavery for domestic and court purposes was widespread throughout Africa. Plantation slavery also occurred, primarily on the eastern coast of Africa and in parts of West Africa. The importance of domestic plantation slavery increased during the 19th century, due to the abolition of the Atlantic slave trade. Many African states dependent on the international slave trade reoriented their economies towards legitimate commerce worked by slave labor.[3]
8
+
9
+ Multiple forms of slavery and servitude have existed throughout African history, and were shaped by indigenous practices of slavery as well as the Roman institution of slavery[citation needed] (and the later Christian views on slavery), the Islamic institutions of slavery via the Arab slave trade, and eventually the Atlantic slave trade.[1]
10
+ Slavery was a part of the economic structure of African societies for many centuries, although the extent varied.[1]
11
+ Ibn Battuta, who visited the ancient kingdom of Mali in the mid-14th century, recounts that the local inhabitants vied with each other in the number of slaves and servants they had, and was himself given a slave boy as a "hospitality gift."[4] In sub-Saharan Africa, the slave relationships were often complex, with rights and freedoms given to individuals held in slavery and restrictions on sale and treatment by their masters.[5] Many communities had hierarchies between different types of slaves: for example, differentiating between those who had been born into slavery and those who had been captured through war.[6]
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+
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+ Travels in the Interior of Africa, Mungo Park, Travels in the Interior of Africa v. II, Chapter XXII – War and Slavery.
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+
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+ The forms of slavery in Africa were closely related to kinship structures. In many African communities, where land could not be owned, enslavement of individuals was used as a means to increase the influence a person had and expand connections.[7] This made slaves a permanent part of a master's lineage, and the children of slaves could become closely connected with the larger family ties.[1] Children of slaves born into families could be integrated into the master's kinship group and rise to prominent positions within society, even to the level of chief in some instances.[6] However, stigma often remained attached, and there could be strict separations between slave members of a kinship group and those related to the master.[7]
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+ Chattel slavery is a specific servitude relationship where the slave is treated as the property of the owner. As such, the owner is free to sell, trade, or treat the slave as he would other pieces of property, and the children of the slave often are retained as the property of the master.[8] There is evidence of long histories of chattel slavery in the Nile River valley and North Africa, but evidence is incomplete about the extent and practices of chattel slavery throughout much of the rest of the continent prior to written records by Arab or European traders.[8][9]
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+
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+ Many slave relationships in Africa revolved around domestic slavery, where slaves would work primarily in the house of the master, but retain some freedoms. Domestic slaves could be considered part of the master's household and would not be sold to others without extreme cause. The slaves could own the profits from their labour (whether in land or in products), and could marry and pass the land on to their children in many cases.[6][10]
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+
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+ Pawnship, or debt bondage slavery, involves the use of people as collateral to secure the repayment of debt. Slave labor is performed by the debtor, or a relative of the debtor (usually a child). Pawnship was a common form of collateral in West Africa. It involved the pledge of a person or a member of that person's family, to serve another person providing credit. Pawnship was related to, yet distinct from, slavery in most conceptualizations, because the arrangement could include limited, specific terms of service to be provided, and because kinship ties would protect the person from being sold into slavery. Pawnship was a common practice throughout West Africa prior to European contact, including among the Akan people, the Ewe people, the Ga people, the Yoruba people, and the Edo people (in modified forms, it also existed among the Efik people, the Igbo people, the Ijaw people, and the Fon people).[11][12]
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+
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+ Military slavery involved the acquisition and training of conscripted military units which would retain the identity of military slaves even after their service.[13] Slave soldier groups would be run by a Patron, who could be the head of a government or an independent warlord, and who would send his troops out for money and his own political interests.[13]
24
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+ This was most significant in the Nile valley (primarily in Sudan and Uganda), with slave military units organized by various Islamic authorities,[13] and with the war chiefs of Western Africa.[14] The military units in Sudan were formed in the 1800s through large-scale military raiding in the area which is currently the countries of Sudan and South Sudan.[13]
26
+
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+ Moreover, a considerable number of the men born between 1800 and 1849 in West African regions (today Ghana and Burkina Faso) and abducted as slaves to serve in the army in Dutch Indonesia. Interestingly, soldiers were on average 3 cm taller than other West African population. Furthermore, data showed, West Africans were shorter than North Europeans but of almost equal height to South Europeans. This was mainly related to the quality of the nutrition and healthcare.[15]
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+
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+ Human sacrifice was common in West African states up to and during the 19th century. Although archaeological evidence is not clear on the issue prior to European contact, in those societies that practiced human sacrifice, slaves became the most prominent victims.[1]
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+ The Annual Customs of Dahomey were the most notorious example of the human sacrifice of slaves, where 500 prisoners would be sacrificed. Sacrifices were carried out all along the West African coast and further inland. Sacrifices were common in the Benin Empire, in what is now Ghana, and in the small independent states in what is now southern Nigeria. In the Ashanti Region, human sacrifice was often combined with capital punishment.[16][17][18]
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+
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+ Many nations such as the Ashanti of present-day Ghana and the Yoruba of present-day Nigeria were involved in slave-trading. Groups such as the Imbangala of Angola and the Nyamwezi of Tanzania would serve as intermediaries or roving bands, waging war on African states to capture people for export as slaves. Historians John Thornton and Linda Heywood of Boston University have estimated that of the Africans captured and then sold as slaves to the New World in the Atlantic slave trade, around 90% were enslaved by fellow Africans who sold them to European traders.[19] Henry Louis Gates, the Harvard Chair of African and African American Studies, has stated that "without complex business partnerships between African elites and European traders and commercial agents, the slave trade to the New World would have been impossible, at least on the scale it occurred."[19]
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+
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+ The entire Bubi ethnic group descends from escaped intertribal slaves owned by various ancient West-central African ethnic groups.[citation needed]
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+ Like most other regions of the world, slavery and forced labor existed in many kingdoms and societies of Africa for hundreds of years.[20][5] According to Ugo Kwokeji, early European reports of slavery throughout Africa in the 1600s are unreliable because they often conflated various forms of servitude as equal to chattel slavery.[21]
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+
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+ The best evidence of slave practices in Africa come from the major kingdoms, particularly along the coast, and there is little evidence of widespread slavery practices in stateless societies.[1][5][6] Slave trading was mostly secondary to other trade relationships; however, there is evidence of a trans-Saharan slave trade route from Roman times which persisted in the area after the fall of the Roman Empire.[8] However, kinship structures and rights provided to slaves (except those captured in war) appears to have limited the scope of slave trading before the start of the Arab slave trade and the Atlantic slave trade.[5]
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+
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+ Slavery in northern Africa dates back to ancient Egypt. The New Kingdom (1558–1080 BC) brought in large numbers of slaves as prisoners of war up the Nile valley and used then for domestic and supervised labor.[23] Ptolemaic Egypt (305 BC–30 BC) used both land and sea routes to bring slaves in.[24]
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+
43
+ Chattel slavery had been legal and widespread throughout North Africa when the region was controlled by the Roman Empire (145 BC – ca. 430 AD), and by the Eastern Romans from 533 to 695). A slave trade bringing Saharans through the desert to North Africa, which existed in Roman times, continued and documentary evidence in the Nile Valley shows it to have been regulated there by treaty.[8] As the Roman republic expanded, it enslaved defeated enemies and Roman conquests in Africa were no exception. For example, Orosius records that Rome enslaved 27,000 people from North Africa in 256 BC.[25] Piracy became an important source of slaves for the Roman Empire and in the 5th century AD pirates would raid coastal North African villages and enslave the captured.[26] Chattel slavery persisted after the fall of the Roman Empire in the largely Christian communities of the region. After the Islamic expansion into most of the region because of the trade expansion across the Sahara,[27] the practices continued and eventually, the assimilative form of slavery spread to major societies on the southern end of the Sahara (such as Mali, Songhai, and Ghana).[1]
44
+ The medieval slave trade in Europe was mainly to the East and South: the Christian Byzantine Empire and the Muslim World were the destinations, Central and Eastern Europe an important source of slaves.[28] Slavery in medieval Europe was so widespread that the Roman Catholic Church repeatedly prohibited it—or at least the export of Christian slaves to non-Christian lands was prohibited at, for example, the Council of Koblenz in 922, the Council of London in 1102, and the Council of Armagh in 1171. Because of religious constraints, the slave trade was carried out in parts of Europe by Iberian Jews (known as Radhanites) who were able to transfer slaves from pagan Central Europe through Christian Western Europe to Muslim countries in Al-Andalus and Africa.[29]
45
+
46
+ The Mamluks were slave soldiers who converted to Islam and served the Muslim caliphs and the Ayyubid Sultans during the Middle Ages. The first Mamluks served the Abbasid caliphs in 9th century Baghdad. Over time, they became a powerful military caste, and on more than one occasion they seized power for themselves, for example, ruling Egypt from 1250–1517. From 1250 Egypt had been ruled by the Bahri dynasty of Kipchak Turk origin. White enslaved people from the Caucasus served in the army and formed an elite corps of troops, eventually revolting in Egypt to form the Burgi dynasty.[30]
47
+ According to Robert Davis between 1 million and 1.25 million Europeans were captured by Barbary pirates and sold as slaves to North Africa and the Ottoman Empire between the 16th and 19th centuries.[31][32] However, to extrapolate his numbers, Davis assumes the number of European slaves captured by Barbary pirates were constant for a 250-year period, stating:
48
+
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+ "There are no records of how many men, women and children were enslaved, but it is possible to calculate roughly the number of fresh captives that would have been needed to keep populations steady and replace those slaves who died, escaped, were ransomed, or converted to Islam. On this basis, it is thought that around 8,500 new slaves were needed annually to replenish numbers - about 850,000 captives over the century from 1580 to 1680. By extension, for the 250 years between 1530 and 1780, the figure could easily have been as high as 1,250,000."[33]
50
+
51
+ Davis' numbers have been disputed by other historians, such as David Earle, who cautions that the true picture of European slaves is clouded by the fact the corsairs also seized non-Christian whites from eastern Europe and black people from west Africa.[33]
52
+
53
+ In addition, the number of slaves traded was hyperactive, with exaggerated estimates relying on peak years to calculate averages for entire centuries, or millennia. Hence, there were wide fluctuations year-to-year, particularly in the 18th and 19th centuries, given slave imports, and also given the fact that, prior to the 1840s, there are no consistent records. Middle East expert John Wright cautions that modern estimates are based on back-calculations from human observation.[34]
54
+
55
+ Such observations, across the late 1500s and early 1600s observers, estimate that around 35,000 European Christian slaves held throughout this period on the Barbary Coast, across Tripoli, Tunis, but mostly in Algiers. The majority were sailors (particularly those who were English), taken with their ships, but others were fishermen and coastal villagers. However, most of these captives were people from lands close to Africa, particularly Spain and Italy.[35]
56
+
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+ The coastal villages and towns of Italy, Portugal, Spain, and Mediterranean islands were frequently attacked by the pirates, and long stretches of the Italian and Spanish coasts were almost completely abandoned by their inhabitants; after 1600 Barbary pirates occasionally entered the Atlantic and struck as far north as Iceland. The most famous corsairs were the Ottoman Barbarossa ("Redbeard"), and his older brother Oruç, Turgut Reis (known as Dragut in the West), Kurtoğlu (known as Curtogoli in the West), Kemal Reis, Salih Reis, and Koca Murat Reis.[32][36]
58
+
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+ In 1544, Hayreddin Barbarossa captured Ischia, taking 4,000 prisoners in the process, and deported to slavery some 9,000 inhabitants of Lipari, almost the entire population.[37] In 1551, Dragut enslaved the entire population of the Maltese island Gozo, between 5,000 and 6,000, sending them to Libya. When pirates sacked Vieste in southern Italy in 1554 they took an estimated 7,000 slaves. In 1555, Turgut Reis sailed to Corsica and ransacked Bastia, taking 6,000 prisoners. In 1558 Barbary corsairs captured the town of Ciutadella, destroyed it, slaughtered the inhabitants, and carried off 3,000 survivors to Istanbul as slaves.[38] In 1563 Turgut Reis landed at the shores of the province of Granada, Spain, and captured the coastal settlements in the area like Almuñécar, along with 4,000 prisoners. Barbary pirates frequently attacked the Balearic islands, resulting in many coastal watchtowers and fortified churches being erected. The threat was so severe that Formentera became uninhabited.[39][40]
60
+
61
+ Early modern sources are full of descriptions of the sufferings of Christian galley slaves of the Barbary corsairs:
62
+
63
+ Those who have not seen a galley at sea, especially in chasing or being chased, cannot well conceive the shock such a spectacle must give to a heart capable of the least tincture of commiseration. To behold ranks and files of half-naked, half-starved, half-tanned meagre wretches, chained to a plank, from whence they remove not for months together (commonly half a year), urged on, even beyond human strength, with cruel and repeated blows on their bare flesh...[41]
64
+
65
+ As late as 1798, the islet near Sardinia was attacked by the Tunisians and over 900 inhabitants were taken away as slaves.
66
+
67
+ Sahrawi-Moorish society in Northwest Africa was traditionally (and still is, to some extent) stratified into several tribal castes, with the Hassane warrior tribes ruling and extracting tribute – horma – from the subservient Berber-descended znaga tribes. Below them ranked servile groups known as Haratin, a black population.[42]
68
+
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+ Enslaved Sub-Saharan Africans were also transported across North Africa into Arabia to do agricultural work because of their resistance to malaria that plagued the Arabia and North Africa at the time of early enslavement. Sub-Saharan Africans were able to endure the malaria-infested lands they were transported to, which is why North Africans were not transported despite their close proximity to Arabia and its surrounding lands.[43]
70
+
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+ In the Horn of Africa, the Christian kings of the Ethiopian Empire often exported pagan Nilotic slaves from their western borderlands, or from newly conquered or reconquered lowland territories.[44] The Somali and Afar Muslim sultanates, such as the medieval Adal Sultanate, through their ports also traded Zanj (Bantu) slaves who were captured from the hinterland.[45]
72
+
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+ Slavery, as practiced in Ethiopia, was essentially domestic and was geared more towards women; this was the trend for most of Africa as well. Women were transported across the Sahara, the Middle East, and the Mediterranean and the Indian Ocean trade more than men.[46] Enslaved people served in the houses of their masters or mistresses, and were not employed to any significant extent for productive purpose. The enslaved were regarded as second-class members of their owners' family.[47] The first attempt to abolish slavery in Ethiopia was made by Emperor Tewodros II (r. 1855–68),[48] although the slave trade was not abolished legally until 1923 with Ethiopia's ascension to the League of Nations.[49] Anti-Slavery Society estimated there were 2 million slaves in the early 1930s, out of an estimated population of between 8 and 16 million.[50] Slavery continued in Ethiopia until the Italian invasion in October 1935, when the institution was abolished by order of the Italian occupying forces.[51] In response to pressure by Western Allies of World War II, Ethiopia officially abolished slavery and involuntary servitude after having regained its independence in 1942.[52][53] On 26 August 1942, Haile Selassie issued a proclamation outlawing slavery.[54]
74
+
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+ In Somali territories, slaves were purchased in the slave market exclusively to do work on plantation grounds.[55] In terms of legal considerations, the customs regarding the treatment of Bantu slaves were established by the decree of Sultans and local administrative legates[disambiguation needed]. Additionally, freedom for these plantation slaves was also often acquired through eventual emancipation, escape, and ransom.[55]
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+
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+ Oral tradition recounts slavery existing in the Kingdom of Kongo from the time of its formation with Lukeni lua Nimi enslaving the Mwene Kabunga whom he conquered to establish the kingdom.[56] Early Portuguese writings show that the Kingdom did have slavery before contact, but that they were primarily war captives from the Kingdom of Ndongo.[56][57]
78
+
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+ Slavery was common along the Upper Congo River, and in the second half of the 18th century the region became a major source of slaves for the Atlantic Slave Trade, when high slave prices on the coast made long-distance slave trading profitable. When the Atlantic trade came to an end, the prices of slaves dropped dramatically, and the regional slave trade grew, dominated by Bobangi traders. The Bobangi also purchased a large number of slaves with profits from selling ivory, who they used to populate their villages. A distinction was made between two different types of slaves in this region; slaves who had been sold by their kin group, typically as a result of undesirable behavior such as adultery, were unlikely to attempt to flee. In addition to those considered socially undesirable, the sale of children was also common in times of famine. Slaves who were captured, however, were likely to attempt to escape and had to be moved hundreds of kilometers from their homes as a safeguard against this.[58]
80
+
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+ The slave trade had a profound impact on this region of Central Africa, completely reshaping various aspects of society. For instance, the slave trade helped to create a robust regional trade network for the foodstuffs and crafted goods of small producers along the river. As the transport of only a few slaves in a canoe was sufficient to cover the cost of a trip and still make a profit, traders could fill any unused space on their canoes with other goods and transport them long distances without a significant markup on price. While the large profits from the Congo River slave trade only went to a small number of traders, this aspect of the trade provided some benefit to local producers and consumers.[59]
82
+
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+ Various forms of slavery were practiced in diverse ways in different communities of West Africa prior to European trade.[20] Even though slavery did exist, it was not nearly as prevalent within most West African societies that were not Islamic before the Trans-Atlantic Slave Trade.[60][61] The prerequisites for slave societies to exist weren't present in West Africa prior to the Atlantic slave trade considering the small market sizes and the lack of a division of labor.[60] Most West African societies were formed in Kinship units which would make slavery a rather marginal part of the production process within them.[1] Slaves within Kinship-based societies would have had almost the same roles that free members had.[1] Martin Klein has said that before the Atlantic trade, slaves in Western Sudan “made up a small part of the population, lived within the household, worked alongside free members of the household, and participated in a network of face-to-face links.”[60] With the development of the trans-Saharan slave trade and the economies of gold in the western Sahel, a number of the major states became organized around the slave trade, including the Ghana Empire, the Mali Empire, and Songhai Empire.[62] However, other communities in West Africa largely resisted the slave trade. The Jola refused to participate in the slave trade up into the end of the seventeenth century, and didn't use slave labor within their own communities until the nineteenth century. The Kru and Baga also fought against the slave trade.[63] The Mossi Kingdoms tried to take over key sites in the trans-Saharan trade and, when these efforts failed, the Mossi became defenders against slave raiding by the powerful states of the western Sahel. The Mossi would eventually enter the slave trade in the 1800s with the Atlantic slave trade being the main market.[62]
84
+
85
+ Senegal was a catalyst for slave trade, and from the Homann Heirs map figure shown, shows a starting point for migration and a firm port of trade. The culture of the Gold Coast was based largely on the power that individuals held, rather than the land cultivated by a family. Western Africa, and specifically places like Senegal, were able to arrive at the development of slavery through analyzing the aristocratic advantages of slavery and what would best suit the region. This sort of governing that used "political tool" of discerning the different labors and methods of assimilative slavery. The domestic and agricultural labor became more evidently primary in Western Africa due to slaves being regarded as these "political tools" of access and status. Slaves often had more wives than their owners, and this boosted the class of their owners. Slaves were not all used for the same purpose. European colonizing countries were participating in the trade to suit the economic needs of their countries. The parallel of "Moorish" traders found in the desert compared to the Portuguese traders that were not as established pointed out the differences in uses of slaves at this point, and where they were headed in the trade.
86
+
87
+ Historian Walter Rodney identified no slavery or significant domestic servitude in early European accounts on the Upper Guinea region[6] and I. A. Akinjogbin contends that European accounts reveal that the slave trade was not a major activity along the coast controlled by the Yoruba people and Aja people before Europeans arrived.[64] In a paper read to the Ethnological Society of London in 1866, the viceroy of Lokoja Mr T. Valentine Robins, who in 1864 accompanied an expedition up the River Niger aboard HMS Investigator, described slavery in the region:
88
+
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+ Upon slavery Mr Robins remarked that it was not what people in England thought it to be. It means, as continually found in this part of Africa, belonging to a family group-there is no compulsory labour, the owner and the slave work together, eat like food, wear like clothing and sleep in the same huts. Some slaves have more wives than their masters. It gives protection to the slaves and everything necessary for their subsistence - food and clothing. A free man is worse off than a slave; he cannot claim his food from anyone.[65]
90
+
91
+ With the beginning of the Atlantic slave trade, demand for slavery in West Africa increased and a number of states became centered on the slave trade and domestic slavery increased dramatically.[66] Hugh Clapperton in 1824 believed that half the population of Kano were enslaved people.[67]
92
+
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+ In the Senegambia region, between 1300 and 1900, close to one-third of the population was enslaved. In early Islamic states of the western Sahel, including Ghana (750–1076), Mali (1235–1645), Segou (1712–1861), and Songhai (1275–1591), about a third of the population were enslaved. In Sierra Leone in the 19th century about half of the population consisted of enslaved people. Among the Vai people, during the 19th century, three quarters of people were slaves. In the 19th century at least half the population was enslaved among the Duala of the Cameroon and other peoples of the lower Niger, the Kongo, and the Kasanje kingdom and Chokwe of Angola. Among the Ashanti and Yoruba a third of the population consisted of enslaved people. The population of the Kanem (1600–1800) was about one-third enslaved. It was perhaps 40% in Bornu (1580–1890). Between 1750 and 1900 from one- to two-thirds of the entire population of the Fulani jihad states consisted of enslaved people. The population of the Sokoto caliphate formed by Hausas in the northern Nigeria and Cameroon was half-enslaved in the 19th century. Slavery was widespread among Taureg peoples and lasted until at least 1975. Among the Adrar 15 percent of people were enslaved, and 75 percent of the Gurma were enslaved.[68]
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+
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+ When British rule was first imposed on the Sokoto Caliphate and the surrounding areas in northern Nigeria at the turn of the 20th century, approximately 2 million to 2.5 million people there were enslaved.[69] Slavery in northern Nigeria was finally outlawed in 1936.[70]
96
+
97
+ With sea trade from the eastern African Great Lakes region to Persia, China, and India during the first millennium AD, slaves are mentioned as a commodity of secondary importance to gold and ivory. When mentioned, the slave trade appears to be of a small-scale and mostly involves slave raiding of women and children along the islands of Kilwa Kisiwani, Madagascar, and Pemba. In places such as Uganda, the experience for women in slavery was different than that of customary slavery practices at the time. The roles assumed were based off gender and position within the society. First one must make the distinction in Ugandan slavery of peasants and slaves. Researchers Shane Doyle and Henri Médard assert the distinction with the following:
98
+
99
+ "Peasants were rewarded for valour in battle by the present of slaves by the lord or chief for whom they had fought. They could be given slaves by relatives who had been promoted to the rank of chiefs, and they could inherit slaves from their fathers. There were the abanyage (those pillaged or stolen in war) as well as the abagule (those bought). All these came under the category of abenvumu or true slaves, that is to say people not free in any sense. In a superior position were the young Ganda given by their maternal uncles into slavery (or pawnship), usually in lieu of debts... Besides such slaves both chiefs and king were served by sons of well to do men who wanted to please them and attract favour for themselves or their children. These were the abasige and formed a big addition to a noble household.... All these different classes of dependents in a household were classed as Medard & Doyle abaddu (male servants) or abazana (female servants) whether they were slave or free-born.(175)"
100
+
101
+ In the Great Lakes region of Africa (around present-day Uganda), linguistic evidence shows the existence of slavery through war capture, trade, and pawning going back hundreds of years; however, these forms, particularly pawning, appear to have increased significantly in the 18th and 19th centuries.[71] These slaves were considered to be more trustworthy than those from the Gold Coast. They were regarded with more prestige because of the training they responded to.
102
+
103
+ The language for slaves in the Great Lakes region varied. This region of water made it easy for capture of slaves and transport. Captive, refugee, slave, peasant were all used in order to describe those in the trade. The distinction was made by where and for what purpose they would be utilized for. Methods like pillage, plunder, and capture were all semantics common in this region to depict the trade.
104
+
105
+ Historians Campbell and Alpers argue that there were a host of different categories of labor in Southeast Africa and that the distinction between slave and free individuals was not particularly relevant in most societies.[72] However, with increasing international trade in the 18th and 19th century, Southeast Africa began to be involved significantly in the Atlantic slave trade; for example, with the king of Kilwa island signing a treaty with a French merchant in 1776 for the delivery of 1,000 slaves per year.[73]
106
+
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+ At about the same time, merchants from Oman, India, and Southeast Africa began establishing plantations along the coasts and on the islands,[74] To provide workers on these plantations, slave raiding and slave holding became increasingly important in the region and slave traders (most notably Tippu Tip) became prominent in the political environment of the region.[73] The Southeast African trade reached its height in the early decades of the 1800s with up to 30,000 slaves sold per year. However, slavery never became a significant part of the domestic economies except in Sultanate of Zanzibar where plantations and agricultural slavery were maintained.[66] Author and historian Timothy Insoll wrote: "Figures record the exporting of 718,000 slaves from the Swahili coast during the 19th century, and the retention of 769,000 on the coast."[75] At various times, between 65 and 90 percent of Zanzibar was enslaved. Along the Kenya coast, 90 percent of the population was enslaved, while half of Madagascar's population was enslaved.[76]
108
+
109
+ Slave relationships in Africa have been transformed through three large-scale processes: the Arab slave trade, the Atlantic slave trade, and the slave emancipation policies and movements in the 19th and 20th centuries. Each of these processes significantly changed the forms, level, and economics of slavery in Africa.[1]
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+
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+ Slave practices in Africa were used during different periods to justify specific forms of European engagement with the peoples of Africa. Eighteenth century writers in Europe claimed that slavery in Africa was quite brutal in order to justify the Atlantic slave trade. Later writers used similar arguments to justify intervention and eventual colonization by European powers to end slavery in Africa.[77]
112
+
113
+ Africans knew of the harsh slavery that awaited slaves in the New World. Many elite Africans visited Europe on slave ships following the prevailing winds through the New World. One example of this occurred when Antonio Manuel, Kongo’s ambassador to the Vatican, went to Europe in 1604, stopping first in Bahia, Brazil, where he arranged to free a countryman who had been wrongfully enslaved. African monarchs also sent their children along these same slave routes to be educated in Europe, and thousands of former slaves eventually returned to settle Liberia and Sierra Leone.[19]
114
+
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+ The enslavement of Africans for eastern markets started before 7th century but remained at low levels until 1750.[78] The trade volume peaked around 1850 but would largely have ended around 1900.[78] Arab slave trade starting in the eighth and ninth centuries AD, began with small-scale movement of people largely from the eastern Great Lakes region and the Sahel. Islamic law allowed slavery, but prohibited slavery involving other pre-existing Muslims; as a result, the main target for slavery were the people who lived in the frontier areas of Islam in Africa.[8] The trade of slaves across the Sahara and across the Indian Ocean also has a long history beginning with the control of sea routes by Afro-Arab traders in the ninth century. It is estimated that, at that time, a few thousand enslaved people were taken each year from the Red Sea and Indian Ocean coast. They were sold throughout the Middle East. This trade accelerated as superior ships led to more trade and greater demand for labour on plantations in the region. Eventually, tens of thousands per year were being taken.[79] On the Swahili Coast, the Afro-Arab slavers captured Bantu peoples from the interior and brought them to the littoral.[80][81] There, the slaves gradually assimilated in the rural areas, particularly on the Unguja and Pemba islands.[80]
116
+
117
+ This changed the slave relationships by creating new forms of employment by slaves (as eunuchs to guard harems, and in military units) and creating conditions for freedom (namely conversion—although it would only free a slave's children).[1][13] Although the level of the trade remained relatively small, the size of total slaves traded grew to a large number over the multiple centuries of its existence.[1] Because of its small and gradual nature, the impact on slavery practices in communities that did not convert to Islam was relatively small.[1] However, in the 1800s, the slave trade from Africa to the Islamic countries picked up significantly. When the European slave trade ended around the 1850s, the slave trade to the east picked up significantly only to be ended with European colonization of Africa around 1900.[66] Between 1500 and 1900, up to 17 million Africans slaves were transported by Muslim traders to the coast of the Indian Ocean, the Middle East, and North Africa.[82]
118
+
119
+ In 1814, Swiss explorer Johann Burckhardt wrote of his travels in Egypt and Nubia, where he saw the practice of slave trading: "I frequently witnessed scenes of the most shameless indecency, which the traders, who were the principal actors, only laughed at. I may venture to state, that very few female slaves who have passed their tenth year, reach Egypt or Arabia in a state of virginity."[83]
120
+
121
+ David Livingstone while talking about the slave trade in East Africa in his journals:
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+
123
+ To overdraw its evil is a simple impossibility.[84]:442
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+
125
+ Livingstone wrote about a group of slaves forced to march by Arab slave traders in the African Great Lakes region when he was travelling there in 1866:
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+
127
+ 19th June 1866 - We passed a woman tied by the neck to a tree and dead, the people of the country explained that she had bene unable to keep up with the other slaves in a gang, and her master had determined that she should not become anyone's property if she recovered.[84]:5626th June 1866 - ... We passed a slave woman shot or stabbed through the body and lying on the path: a group of mon stood about a hundred yards off on one side, and another of the women on the other side, looking on; they said an Arab who passed early that morning had done it in anger at losing the price he had given for her, because she was unable to walk any longer.
128
+ 27th June 1866 - To-day we came upon a man dead from starvation, as he was very thin. One of our men wandered and found many slaves with slave-sticks on, abandoned by their masters from want of food; they were too weak to be able to speak or say where they had come from; some were quite young.[84]:62
129
+
130
+ Zanzibar was once East Africa's main slave-trading port, and under Omani Arabs in the 19th century as many as 50,000 slaves were passing through the city each year.[85]
131
+
132
+ The Atlantic slave trade or transatlantic slave trade took place across the Atlantic Ocean from the 15th through to the 19th centuries. According to Patrick Manning, the Atlantic slave trade was significant in transforming Africans from a minority of the global population of slaves in 1600 into the overwhelming majority by 1800 and by 1850 the number of African slaves within Africa exceeded those in the Americas.[86]
133
+
134
+ The slave trade was transformed from a marginal aspect of the economies into the largest sector in a relatively short span. In addition, agricultural plantations increased significantly and became a key aspect in many societies.[1] Economic urban centers that served as the root of main trade routes shifted towards the West coast.[87] At the same time, many African communities relocated far away from slave trade routes, often protecting themselves from the Atlantic slave trade but hindering economic and technological development at the same time.[88]
135
+
136
+ In many African societies traditional lineage slavery became more like chattel slavery due to an increased work demand.[89] This resulted in a general decrease in quality of life, working conditions, and status of slaves in West African societies. Assimilative slavery was increasingly replaced with chattel slavery. Assimilitave slavery in Africa often allowed eventual freedom and also significant cultural, social, and/or economic influence. Slaves were often treated as part of their owner's family, rather than simply property.[89]
137
+
138
+ The distribution of gender among enslaved peoples under traditional lineage slavery saw women as more desirable slaves due to demands for domestic labor and for reproductive reasons.[89] Male slaves were used for more physical agricultural labor,[90] but as more enslaved men were taken to the West Coast and across the Atlantic to the New World, female slaves were increasingly used for physical and agricultural labour and polygyny also increased. Chattel slavery in America was highly demanding because of the physical nature of plantation work and this was the most common destination for male slaves in the New World.[89]
139
+
140
+ It has been argued that a decrease in able-bodied people as a result of the Atlantic slave trade limited many societies ability to cultivate land and develop. Many scholars argue that the transatlantic slave trade, left Africa underdeveloped, demographically unbalanced, and vulnerable to future European colonization.[88]
141
+
142
+ The first Europeans to arrive on the coast of Guinea were the Portuguese; the first European to actually buy enslaved Africans in the region of Guinea was Antão Gonçalves, a Portuguese explorer in 1441 AD. Originally interested in trading mainly for gold and spices, they set up colonies on the uninhabited islands of São Tomé. In the 16th century the Portuguese settlers found that these volcanic islands were ideal for growing sugar. Sugar growing is a labour-intensive undertaking and Portuguese settlers were difficult to attract due to the heat, lack of infrastructure, and hard life. To cultivate the sugar the Portuguese turned to large numbers of enslaved Africans. Elmina Castle on the Gold Coast, originally built by African labour for the Portuguese in 1482 to control the gold trade, became an important depot for slaves that were to be transported to the New World.[91]
143
+
144
+ The Spanish were the first Europeans to use enslaved Africans in America on islands such as Cuba and Hispaniola,[92] where the alarming death rate in the native population had spurred the first royal laws protecting the native population (Laws of Burgos, 1512–13). The first enslaved Africans arrived in Hispaniola in 1501 soon after the Papal Bull of 1493 gave almost all of the New World to Spain.[93]
145
+
146
+ In Igboland, for example, the Aro oracle (the Igbo religious authority) began condemning more people to slavery due to small infractions that previously probably wouldn't have been punishable by slavery, thus increasing the number of enslaved men available for purchase.[89]
147
+
148
+ The Atlantic slave trade peaked in the late 18th century, when the largest number of people were bought or captured from West Africa and taken to the Americas.[94] The increase of demand for slaves due to the expansion of European colonial powers to the New World made the slave trade much more lucrative to the West African powers, leading to the establishment of a number of actual West African empires thriving on slave trade.
149
+ These included the Oyo empire (Yoruba), Kong Empire, Imamate of Futa Jallon, Imamate of Futa Toro, Kingdom of Koya, Kingdom of Khasso, Kingdom of Kaabu, Fante Confederacy, Ashanti Confederacy, and the kingdom of Dahomey. These kingdoms relied on a militaristic culture of constant warfare to generate the great numbers of human captives required for trade with the Europeans.[95] It is documented in the Slave Trade Debates of England in the early 19th century: "All the old writers concur in stating not only that wars are entered into for the sole purpose of making slaves, but that they are fomented by Europeans, with a view to that object."[96] The gradual abolition of slavery in European colonial empires during the 19th century again led to the decline and collapse of these African empires. When European powers began to stop the Atlantic slave trade, this caused a further change in that large holders of slaves in Africa began to exploit enslaved people on plantations and other agricultural products.[97]
150
+
151
+ The final major transformation of slave relationships came with the inconsistent emancipation efforts starting in the mid-19th century. As European authorities began to take over large parts of inland Africa starting in the 1870s, the colonial policies were often confusing on the issue. For example, even when slavery was deemed illegal, colonial authorities would return escaped slaves to their masters.[1] Slavery persisted in some countries under colonial rule, and in some instances it was not until independence that slavery practices were significantly transformed.[98] Anti-colonial struggles in Africa often brought slaves and former slaves together with masters and former masters to fight for independence; however, this cooperation was short-lived and following independence political parties would often form based upon the stratifications of slaves and masters.[66]
152
+
153
+ In some parts of Africa, slavery and slavery-like practices continue to this day, particularly the illegal trafficking of women and children.[99] The problem has proven to be difficult for governments and civil society to eliminate.[100]
154
+
155
+ Efforts by Europeans against slavery and the slave trade began in the late 18th century and had a large impact on slavery in Africa. Portugal was the first country in the continent to abolish slavery in metropolitan Portugal and Portuguese India by a bill issued on 12 February 1761, but this did not affect their colonies of Brazil and Portuguese Africa. France abolished slavery in 1794. However, slavery was again allowed by Napoleon in 1802 and not abolished for good until 1848. In 1803, Denmark-Norway became the first country from Europe to implement a ban on the slave trade. Slavery itself was not banned until 1848.[101] Britain followed in 1807 with the passage of the Abolition of the Slave Trade Act by Parliament. This law allowed stiff fines, increasing with the number of slaves transported, for captains of slave ships. Britain followed this with the Slavery Abolition Act 1833 which freed all slaves in the British Empire. British pressure on other countries resulted in them agreeing to end the slave trade from Africa. For example, the 1820 U.S. Law on Slave Trade made slave trading piracy, punishable by death.[102] In addition, the Ottoman Empire abolished slave trade from Africa in 1847 under British pressure.[103]
156
+
157
+ By 1850, the year that the last major Atlantic slave trade participant (Brazil) passed the Eusébio de Queirós Law banning the slave trade,[104] the slave trades had been significantly slowed and in general only illegal trade went on. Brazil continued the practice of slavery and was a major source for illegal trade until about 1870 and the abolition of slavery became permanent in 1888 when Princess Isabel of Brazil and Minister Rodrigo Silva (son-in-law of senator Eusebio de Queiroz) banned the practice.[66] The British took an active approach to stopping the illegal Atlantic slave trade during this period. The West Africa Squadron was credited with capturing 1,600 slave ships between 1808 and 1860, and freeing 150,000 Africans who were aboard these ships.[105] Action was also taken against African leaders who refused to agree to British treaties to outlaw the trade, for example against ‘the usurping King of Lagos’, deposed in 1851. Anti-slavery treaties were signed with over 50 African rulers.
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+ According to Patrick Manning, internal slavery was most important to Africa in the second half of the 19th century, stating "if there is any time when one can speak of African societies being organized around a slave mode production, [1850–1900] was it". The abolition of the Atlantic slave trade resulted in the economies of African states dependent on the trade being reorganized towards domestic plantation slavery and legitimate commerce worked by slave labor. Slavery before this period was generally domestic.[66][3]
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+ The continuing anti-slavery movement in Europe became an excuse and a casus belli for the European conquest and colonization of much of the African continent.[77] It was the central theme of the Brussels Anti-Slavery Conference 1889-90. In the late 19th century, the Scramble for Africa saw the continent rapidly divided between imperialistic European powers, and an early but secondary focus of all colonial regimes was the suppression of slavery and the slave trade. Seymour Drescher argues that European interests in abolition were primarily motivated by economic and imperial goals.[106] Despite slavery often being a justification behind conquest, colonial regimes often ignored slavery or allowed slavery practices to continue. This was because the colonial state depended on the cooperation of indigenous political and economic structures which were heavily involved in slavery. As a result, early colonial policies usually sought to end slave trading while regulating existing slave practices and weakening the power of slave maaters.[61] Furthermore the early colonial states had weak effective control over their territories, which precluded efforts to widespread abolition. Abolition attempts became more concrete later during the colonial period.[61]
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+ There were many causes for the decline and abolition of slavery in Africa during the colonial period including colonial abolition policies, various economic changes, and slave resistance. The economic changes during the colonial period, including the rise of wage labor and cash crops, hastened the decline of slavery by offering new economic opportunities to slaves. The abolition of slave raiding and the end of wars between African states drastically reduced the supply of slaves. Slaves would take advantage of early colonial laws that nominally abolished slavery and would migrate away from their masters although these laws often were intended to regulate slavery more than actually abolish it. This migration led to more concrete abolition efforts by colonial governments.[61][107][1][61]
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+ Following conquest and abolition by the French, over a million slaves in French West Africa fled from their masters to earlier homes between 1906 and 1911.[108] In Madagascar over 500,000 slaves were freed following French abolition in 1896.[109] In response to this pressure, Ethiopia officially abolished slavery in 1932, Sokoto Caliphate abolished slavery in 1900, and the rest of the Sahel in 1911. Colonial nations were mostly successful in this aim, though slavery is still very active in Africa even though it has gradually moved to a wage economy. Independent nations attempting to westernize or impress Europe sometimes cultivated an image of slavery suppression, even as they, in the case of Egypt, hired European soldiers like Samuel White Baker's expedition up the Nile. Slavery has never been eradicated in Africa, and it commonly appears in African states, such as Chad, Ethiopia, Mali, Niger, and Sudan, in places where law and order have collapsed.[110]
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+ Although outlawed in all countries today, slavery is practiced in secret in many parts of the world.[111] There are an estimated 30 million victims of slavery worldwide.[112] In Mauritania alone, up to 600,000 men, women and children, or 20% of the population, are enslaved, many of them used as bonded labour.[113][114] Slavery in Mauritania was finally criminalized in August 2007.[115] During the Second Sudanese Civil War people were taken into slavery; estimates of abductions range from 14,000 to 200,000.[116] In Niger, where the practice of slavery was outlawed in 2003, a study found that almost 8% of the population are still slaves.[117][118]
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+ Slavery and the slave trades had a significant impact on the size of the population and the gender distribution throughout much of Africa. The precise impact of these demographic shifts has been an issue of significant debate.[119] The Atlantic slave trade took 70,000 people, primarily from the west coast of Africa, per year at its peak in the mid-1700s.[66] The Arab slave trade involved the capture of peoples from the continental interior, who were then shipped overseas through ports on the Red Sea and elsewhere.[120] It peaked at 10,000 people bartered per year in the 1600s.[66] According to Patrick Manning, there was a consistent population decrease in large parts of Sub-Saharan Africa as a result of these slave trades. This population decline throughout West Africa from 1650 until 1850 was exacerbated by the preference of slave traders for male slaves. It is important to note that this preference only existed in the transatlantic slave trade. More female slaves than male were traded across the continent of Africa.[46][66] In eastern Africa, the slave trade was multi-directional and changed over time. To meet the demand for menial labor, Zanj slaves captured from the southern interior were sold through ports on the northern seaboard in cumulatively large numbers over the centuries to customers in the Nile Valley, Horn of Africa, Arabian Peninsula, Persian Gulf, India, Far East and the Indian Ocean islands.[120]
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+ The extent of slavery within Africa and the trade in slaves to other regions is not known precisely. Although the Atlantic slave trade has been best studied, estimates range from 8 million people to 20 million.[121] The Trans-Atlantic Slave Trade Database estimates that the Atlantic slave trade took around 12.8 million people between 1450 and 1900.[1][122] The slave trade across the Sahara and Red Sea from the Sahara, the Horn of Africa, and East Africa, has been estimated at 6.2 million people between 600 and 1600.[1] Although the rate decreased from East Africa in the 1700s, it increased in the 1800s and is estimated at 1.65 million for that century.[1]
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+ Estimates by Patrick Manning are that about 12 million slaves entered the Atlantic trade between the 16th and 19th century, but about 1.5 million died on board ship.[123] About 10.5 million slaves arrived in the Americas.[123] Besides the slaves who died on the Middle Passage, more Africans likely died during the slave raids in Africa and forced marches to ports. Manning estimates that 4 million died inside Africa after capture, and many more died young.[123] Manning's estimate covers the 12 million who were originally destined for the Atlantic, as well as the 6 million destined for Asian slave markets and the 8 million destined for African markets.[123]
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+ The demographic effects of the slave trade are some of the most controversial and debated issues. Walter Rodney argued that the export of so many people had been a demographic disaster and had left Africa permanently disadvantaged when compared to other parts of the world, and that this largely explains that continent's continued poverty.[124] He presents numbers that show that Africa's population stagnated during this period, while that of Europe and Asia grew dramatically. According to Rodney all other areas of the economy were disrupted by the slave trade as the top merchants abandoned traditional industries to pursue slaving and the lower levels of the population were disrupted by the slaving itself.
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+ Others have challenged this view. J. D. Fage compared the number effect on the continent as a whole. David Eltis has compared the numbers to the rate of emigration from Europe during this period. In the 19th century alone over 50 million people left Europe for the Americas, a far higher rate than were ever taken from Africa.[125]
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+ Others in turn challenged that view. Joseph E. Inikori argues the history of the region shows that the effects were still quite deleterious. He argues that the African economic model of the period was very different from the European, and could not sustain such population losses. Population reductions in certain areas also led to widespread problems. Inikori also notes that after the suppression of the slave trade Africa's population almost immediately began to rapidly increase, even prior to the introduction of modern medicines.[126]
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+ There is a longstanding debate among analysts and scholars about the destructive impacts of the slave trades.[20] It is often claimed that the slave trade undermined local economies and political stability as villages' vital labour forces were shipped overseas as slave raids and civil wars became commonplace. With the rise of a large commercial slave trade, driven by European needs, enslaving your enemy became less a consequence of war, and more and more a reason to go to war.[127] The slave trade was claimed to have impeded the formation of larger ethnic groups, causing ethnic factionalism and weakening the formation for stable political structures in many places. It also is claimed to have reduced the mental health and social development of African people.[128]
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+ In contrast to these arguments, J. D. Fage asserts that slavery did not have a wholly disastrous effect on the societies of Africa.[129] Slaves were an expensive commodity, and traders received a great deal in exchange for each enslaved person. At the peak of the slave trade hundreds of thousands of muskets, vast quantities of cloth, gunpowder, and metals were being shipped to Guinea. Most of this money was spent on British-made firearms (of very poor quality) and industrial-grade alcohol. Trade with Europe at the peak of the slave trade—which also included significant exports of gold and ivory—was some 3.5 million pounds Sterling per year. By contrast, the trade of the United Kingdom, the economic superpower of the time, was about 14 million pounds per year over this same period of the late 18th century. As Patrick Manning has pointed out, the vast majority of items traded for slaves were common rather than luxury goods. Textiles, iron ore, currency, and salt were some of the most important commodities imported as a result of the slave trade, and these goods were spread within the entire society raising the general standard of living.[20]
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+ Although debated, it is argued that the Atlantic slave trade devastated the African economy. In 19th century Yoruba Land, economic activity was described to be at its lowest ever while life and property were being taken daily, and normal living was in jeopardy because of the fear of being kidnapped.[130] (Onwumah, Imhonopi, Adetunde,2019)
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+ Karl Marx in his economic history of capitalism, [Das Kapital], claimed that "...the turning of Africa into a warren for the commercial hunting of black-skins [that is, the slave trade], signalled the rosy dawn of the era of capitalist production." He argued that the slave trade was part of what he termed the "primitive accumulation" of European capital, the non-capitalist accumulation of wealth that preceded and created the financial conditions for Britain's industrialisation and the advent of the capitalist mode of production.[131]
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+ Eric Williams has written about the contribution of Africans on the basis of profits from the slave trade and slavery, arguing that the employment of those profits were used to help finance Britain’s industrialisation. He argues that the enslavement of Africans was an essential element to the Industrial Revolution, and that European wealth was, in part, a result of slavery, but that by the time of its abolition it had lost its profitability and it was in Britain's economic interest to ban it.[132] Joseph Inikori has written that the British slave trade was more profitable than the critics of Williams believe. Other researchers and historians have strongly contested what has come to be referred to as the "Williams thesis" in academia: David Richardson has concluded that the profits from the slave trade amounted to less than 1% of domestic investment in Britain,[133] and economic historian Stanley Engerman finds that even without subtracting the associated costs of the slave trade (e.g., shipping costs, slave mortality, mortality of whites in Africa, defense costs) or reinvestment of profits back into the slave trade, the total profits from the slave trade and of West Indian plantations amounted to less than 5% of the British economy during any year of the Industrial Revolution.[134] Historian Richard Pares, in an article written before Williams’ book, dismisses the influence of wealth generated from the West Indian plantations upon the financing of the Industrial Revolution, stating that whatever substantial flow of investment from West Indian profits into industry there was occurred after emancipation, not before.[135] Findlay and O'Rourke noted that the figures presented by O'Brien (1982) to back his claim that "the periphery was peripheral" suggest the opposite, with profits from the periphery 1784–1786 being £5.66 million when there was £10.30 million total gross investment in the British economy and similar proportions for 1824–1826. They note that dismissing the profits of the enslavement of human beings from significance because it was a "small share of national income", could be used to argue that there was no industrial revolution, since modern industry provided only a small share of national income and that it is a mistake to assume that small size is the same as small significance. Findlay and O'Rourke also note that the share of American export commodities produced by enslaved human beings, rose from 54% between 1501 and 1550 to 82.5% between 1761–1780.[136]
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+ Seymour Drescher and Robert Anstey argue the slave trade remained profitable until the end, because of innovations in agriculture, and that moralistic reform, not economic incentive, was primarily responsible for abolition.[137]
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+ A similar debate has taken place about other European nations. The French slave trade, it is argued, was more profitable than alternative domestic investments, and probably encouraged capital accumulation before the Industrial Revolution and Napoleonic Wars.[138]
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+ Maulana Karenga states the effects of the Atlantic slave trade in African captives: "[T]he morally monstrous destruction of human possibility involved redefining African humanity to the world, poisoning past, present and future relations with others who only know us through this stereotyping and thus damaging the truly human relations among people of today". He says that it constituted the destruction of culture, language, religion and human possibility.[139]
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+ A tram (in North America streetcar or trolley) is a rail vehicle that runs on tramway tracks along public urban streets; some include segments of segregated right-of-way.[1][2] The lines or networks operated by tramcars are called tramways. Historically the term electric street railways was also used in the United States. In the United States, the term tram has sometimes been used for rubber-tired trackless trains, which are unrelated to other kinds of trams.
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+ Tram vehicles are usually lighter and shorter than main line and rapid transit trains. Today, most trams use electrical power, usually fed by a pantograph sliding on an overhead line; older systems may use a trolley pole or a bow collector. In some cases, a contact shoe on a third rail is used. If necessary, they may have dual power systems—electricity in city streets and diesel in more rural environments. Occasionally, trams also carry freight.
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+ Trams are now commonly included in the wider term "light rail",[3] which also includes grade-separated systems. Some trams, known as tram-trains, may have segments that run on mainline railway tracks, similar to interurban systems. The differences between these modes of rail transport are often indistinct and a given system may combine multiple features.
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+ One of the advantages over earlier forms of transit was the low rolling resistance of metal wheels on steel rails, allowing the trams to haul a greater load for a given effort. Problems included the high total cost of ownership of horses. Electric trams largely replaced animal power in the late 19th and early 20th centuries. Improvements in other vehicles such as buses led to decline of trams in the mid 20th century. However, trams have seen resurgence in recent years.
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+ The English terms tram and tramway are derived from the Scots word tram,[4] referring respectively to a type of truck (goods wagon or freight railroad car) used in coal mines and the tracks on which they ran. The word tram probably derived from Middle Flemish trame ("beam, handle of a barrow, bar, rung"). The identical word la trame with the meaning "crossbeam" is also used in the French language. Etymologists believe that the word tram refers to the wooden beams the railway tracks were initially made of before the railroad pioneers switched to the much more wear-resistant tracks made of iron and, later, steel.[5] The word Tram-car is attested from 1873.[6]
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+ Although the terms tram and tramway have been adopted by many languages, they are not used universally in English; North Americans prefer streetcar, trolley, or trolleycar. The term streetcar is first recorded in 1840, and originally referred to horsecars. When electrification came, Americans began to speak of trolleycars or later, trolleys. A widely held belief holds the word to derive from the troller (said to derive from the words traveler and roller), a four-wheeled device that was dragged along dual overhead wires by a cable that connected the troller to the top of the car and collected electrical power from the overhead wires;[7] this portmanteau derivation is, however, most likely folk etymology. "Trolley" and variants refer to the verb troll, meaning "roll" and probably derived from Old French,[8] and cognate uses of the word were well established for handcarts and horse drayage, as well as for nautical uses.[9]
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+ The alternative North American term 'trolley' may strictly speaking be considered incorrect, as the term can also be applied to cable cars, or conduit cars that instead draw power from an underground supply. Conventional diesel tourist buses decorated to look like streetcars are sometimes called trolleys in the US (tourist trolley). Furthering confusion, the term tram has instead been applied to open-sided, low-speed segmented vehicles on rubber tires generally used to ferry tourists short distances, for example on the Universal Studios backlot tour and, in many countries, as tourist transport to major destinations. The term may also apply to an aerial ropeway, e.g. the Roosevelt Island Tramway.
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+ Although the use of the term trolley for tram was not adopted in Europe, the term was later associated with the trolleybus, a rubber-tyred vehicle running on hard pavement, which draws its power from pairs of overhead wires. These electric buses, which use twin trolley poles, are also called trackless trolleys (particularly in the northeastern US), or sometimes simply trolleys (in the UK, as well as the Pacific Northwest, including Seattle, and Vancouver).
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+ The New South Wales government in Australia has decided to use the term "light rail" for their trams.
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+ The history of trams, streetcars or trolley systems, began in the early nineteenth century. It can be divided up into several discrete periods defined by the principal means of motive power used.
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+ The world's first passenger train or tram was the Swansea and Mumbles Railway, in Wales, UK. The Mumbles Railway Act was passed by the British Parliament in 1804, and horse-drawn service started in 1807.[10] The service closed in 1827, but was restarted in 1860, again using horses.[11] It was worked by steam from 1877, and then, from 1929, by very large (106-seater) electric tramcars, until closure in 1961.[citation needed] The Swansea and Mumbles Railway was something of a one-off however, and no street tramway would appear in Britain until 1860 when one was built in Birkenhead by the American George Francis Train.[12]
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+ Street railways developed in America before Europe, largely due to the poor paving of the streets in American cities which made them unsuitable for horsebuses, which were then common on the well-paved streets of European cities. Running the horsecars on rails allowed for a much smoother ride. There are records of a street railway running in Baltimore as early as 1828, however the first authenticated streetcar in America, was the New York and Harlem Railroad developed by the Irish coach builder John Stephenson, in New York City which began service in the year 1832.[13][14] The New York and Harlem Railroad's Fourth Avenue Line ran along the Bowery and Fourth Avenue in New York City. It was followed in 1835 by the New Orleans and Carrollton Railroad in New Orleans, Louisiana,[15] which still operates as the St. Charles Streetcar Line. Other American cities did not follow until the 1850s, after which the "animal railway" became an increasingly common feature in the larger towns.[15]
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+ The first permanent tram line in continental Europe was opened in Paris in 1855 by Alphonse Loubat who had previously worked on American streetcar lines.[16] The tram was developed in numerous cities of Europe (some of the most extensive systems were found in Berlin, Budapest, Birmingham, Leningrad, Lisbon, London, Manchester, Paris).
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+ The first tram in South America opened in 1858 in Santiago, Chile. The first trams in Australia opened in 1860 in Sydney. Africa's first tram service started in Alexandria on 8 January 1863. The first trams in Asia opened in 1869 in Batavia (now Jakarta), Netherlands East Indies (now Indonesia).
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+ Problems with horsecars included the fact that any given animal could only work so many hours on a given day, had to be housed, groomed, fed and cared for day in and day out, and produced prodigious amounts of manure, which the streetcar company was charged with storing and then disposing of. Since a typical horse pulled a streetcar for about a dozen miles a day and worked for four or five hours, many systems needed ten or more horses in stable for each horsecar.
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+ Horsecars were largely replaced by electric-powered trams following the improvement of an overhead trolley system on trams for collecting electricity from overhead wires by Frank J. Sprague. His spring-loaded trolley pole used a wheel to travel along the wire. In late 1887 and early 1888, using his trolley system, Sprague installed the first successful large electric street railway system in Richmond, Virginia. Within a year, the economy of electric power had replaced more costly horsecars in many cities. By 1889, 110 electric railways incorporating Sprague's equipment had been begun or planned on several continents.[17]
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+ Horses continued to be used for light shunting well into the 20th century, and many large metropolitan lines lasted into the early 20th century. New York City had a regular horsecar service on the Bleecker Street Line until its closure in 1917.[18] Pittsburgh, Pennsylvania, had its Sarah Street line drawn by horses until 1923. The last regular mule-drawn cars in the US ran in Sulphur Rock, Arkansas, until 1926 and were commemorated by a U.S. postage stamp issued in 1983.[19] The last mule tram service in Mexico City ended in 1932, and a mule tram in Celaya, Mexico, survived until 1954.[20] The last horse-drawn tram to be withdrawn from public service in the UK took passengers from Fintona railway station to Fintona Junction one mile away on the main Omagh to Enniskillen railway in Northern Ireland. The tram made its last journey on 30 September 1957 when the Omagh to Enniskillen line closed. The "van" now lies at the Ulster Transport Museum.
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+ Horse-drawn trams still operate on the 1876-built Douglas Bay Horse Tramway in the Isle of Man, and at the 1894-built horse tram at Victor Harbor in South Australia. New horse-drawn systems have been established at the Hokkaidō Museum in Japan and also in Disneyland. A horse tram route in Polish gmina Mrozy, first built in 1902, was reopened in 2012.
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+
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+ The first mechanical trams were powered by steam. Generally, there were two types of steam tram. The first and most common had a small steam locomotive (called a tram engine in the UK) at the head of a line of one or more carriages, similar to a small train. Systems with such steam trams included Christchurch, New Zealand; Sydney, Australia; other city systems in New South Wales; Munich, Germany (from August 1883 on),[21] British India (Pakistan) (from 1885) and the Dublin & Blessington Steam Tramway (from 1888) in Ireland. Steam tramways also were used on the suburban tramway lines around Milan and Padua; the last Gamba de Legn ("Peg-Leg") tramway ran on the Milan-Magenta-Castano Primo route in late 1957.[22]
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+
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+ The other style of steam tram had the steam engine in the body of the tram, referred to as a tram engine (UK) or steam dummy (US). The most notable system to adopt such trams was in Paris. French-designed steam trams also operated in Rockhampton, in the Australian state of Queensland between 1909 and 1939. Stockholm, Sweden, had a steam tram line at the island of Södermalm between 1887 and 1901.
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+ Tram engines usually had modifications to make them suitable for street running in residential areas. The wheels, and other moving parts of the machinery, were usually enclosed for safety reasons and to make the engines quieter. Measures were often taken to prevent the engines from emitting visible smoke or steam. Usually the engines used coke rather than coal as fuel to avoid emitting smoke; condensers or superheating were used to avoid emitting visible steam. A major drawback of this style of tram was the limited space for the engine, so that these trams were usually underpowered. Steam tram engines faded out around the 1890s to 1900s, being replaced by electric trams.
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+ Another motive system for trams was the cable car, which was pulled along a fixed track by a moving steel cable. The power to move the cable was normally provided at a "powerhouse" site a distance away from the actual vehicle. The London and Blackwall Railway, which opened for passengers in east London, England, in 1840 used such a system.[23]
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+ The first practical cable car line was tested in San Francisco, in 1873. Part of its success is attributed to the development of an effective and reliable cable grip mechanism, to grab and release the moving cable without damage. The second city to operate cable trams was Dunedin in New Zealand, from 1881 to 1957.
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+ The most extensive cable system in the US was built in Chicago, having been built in stages between 1859 and 1892. New York City developed multiple cable car lines, that operated from 1883 to 1909.[24] Los Angeles also had several cable car lines, including the Second Street Cable Railroad, which operated from 1885 to 1889, and the Temple Street Cable Railway, which operated from 1886 to 1898.
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+ From 1885 to 1940, the city of Melbourne, Victoria, Australia operated one of the largest cable systems in the world, at its peak running 592 trams on 75 kilometres (47 mi) of track. There were also two isolated cable lines in Sydney, New South Wales, Australia; the North Sydney line from 1886 to 1900,[25] and the King Street line from 1892 to 1905.
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+
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+ In Dresden, Germany, in 1901 an elevated suspended cable car following the Eugen Langen one-railed floating tram system started operating. Cable cars operated on Highgate Hill in North London and Kennington to Brixton Hill In South London.[when?] They also worked around "Upper Douglas" in the Isle of Man from 1897 to 1929 (cable car 72/73 is the sole survivor of the fleet).
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+
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+ Cable cars suffered from high infrastructure costs, since an expensive system of cables, pulleys, stationary engines and lengthy underground vault structures beneath the rails had to be provided. They also required physical strength and skill to operate, and alert operators to avoid obstructions and other cable cars. The cable had to be disconnected ("dropped") at designated locations to allow the cars to coast by inertia, for example when crossing another cable line. The cable would then have to be "picked up" to resume progress, the whole operation requiring precise timing to avoid damage to the cable and the grip mechanism. Breaks and frays in the cable, which occurred frequently, required the complete cessation of services over a cable route while the cable was repaired. Due to overall wear, the entire length of cable (typically several kilometres) would have to be replaced on a regular schedule. After the development of reliable electrically powered trams, the costly high-maintenance cable car systems were rapidly replaced in most locations.
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+
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+ Cable cars remained especially effective in hilly cities, since their nondriven wheels would not lose traction as they climbed or descended a steep hill. The moving cable would physically pull the car up the hill at a steady pace, unlike a low-powered steam or horse-drawn car. Cable cars do have wheel brakes and track brakes, but the cable also helps restrain the car to going downhill at a constant speed. Performance in steep terrain partially explains the survival of cable cars in San Francisco.
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+ The San Francisco cable cars, though significantly reduced in number, continue to perform a regular transportation function, in addition to being a well-known tourist attraction. A single cable line also survives in Wellington, New Zealand (rebuilt in 1979 as a funicular but still called the "Wellington Cable Car"). Another system, actually two separate cable lines with a shared power station in the middle, operates from the Welsh town of Llandudno up to the top of the Great Orme hill in North Wales, UK.
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+ In the late 19th and early 20th centuries a number of systems in various parts of the world employed trams powered by gas, naphtha gas or coal gas in particular. Gas trams are known to have operated between Alphington and Clifton Hill in the northern suburbs of Melbourne, Australia (1886–1888); in Berlin and Dresden, Germany; in Estonia (1921–1951); between Jelenia Góra, Cieplice, and Sobieszów in Poland (from 1897); and in the UK at Lytham St Annes, Neath (1896–1920), and Trafford Park, Manchester (1897–1908).
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+
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+ On 29 December 1886 the Melbourne newspaper The Argus reprinted a report from the San Francisco Bulletin that Mr Noble had demonstrated a new 'motor car' for tramways 'with success'. The tramcar 'exactly similar in size, shape, and capacity to a cable grip car' had the 'motive power' of gas 'with which the reservoir is to be charged once a day at power stations by means of a rubber hose'. The car also carried an electricity generator for 'lighting up the tram and also for driving the engine on steep grades and effecting a start'.[26]
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+ Comparatively little has been published about gas trams. However, research on the subject was carried out for an article in the October 2011 edition of "The Times", the historical journal of the Australian Association of Timetable Collectors, now the Australian Timetable Association.[27][28][29][30]
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+ A tram system powered by compressed natural gas was due to open in Malaysia in 2012,[31] but the news about the project appears to have dried up.
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+ The world's first electric tram line operated in Sestroretsk near Saint Petersburg, Russia, invented and tested by Fyodor Pirotsky in 1875.[32][33] Later, using a similar technology, Pirotsky put into service the first public electric tramway in St. Petersburg, which operated only during September 1880.[34]
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+ The second demonstrative tramway was presented by Siemens & Halske at the 1879 Berlin Industrial Exposition.
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+ The first public electric tramway used for permanent service was the Gross-Lichterfelde tramway in Lichterfelde near Berlin in Germany, which opened in 1881. It was built by Werner von Siemens who contacted Pirotsky. This was world's first commercially successful electric tram. It initially drew current from the rails, with overhead wire being installed in 1883.[35]
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+ In Britain, Volk's Electric Railway was opened in 1883 in Brighton). This two kilometer line along the seafront, re-gauged to 2 ft 9 in (838 mm) in 1884, remains in service to this day and is the oldest operating electric tramway in the world. Also in 1883, Mödling and Hinterbrühl Tram was opened near Vienna in Austria. It was the first tram in the world in regular service that was run with electricity served by an overhead line with pantograph current collectors. The Blackpool Tramway was opened in Blackpool, UK on 29 September 1885 using conduit collection along Blackpool Promenade. This system is still in operation in a modernised form.[36]
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+
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+ Earliest tram system in Canada was by John Joseph Wright, brother of the famous mining entrepreneur Whitaker Wright, in Toronto in 1883, introducing electric trams in 1892. In the US, multiple functioning experimental electric trams were exhibited at the 1884 World Cotton Centennial World's Fair in New Orleans, Louisiana, but they were not deemed good enough to replace the Lamm fireless engines then propelling the St. Charles Avenue Streetcar in that city. The first commercial installation of an electric streetcar in the United States was built in 1884 in Cleveland, Ohio and operated for a period of one year by the East Cleveland Street Railway Company.[37] Trams were operated in Richmond, Virginia, in 1888, on the Richmond Union Passenger Railway built by Frank J. Sprague. Sprague later developed multiple unit control, first demonstrated in Chicago in 1897, allowing multiple cars to be coupled together and operated by a single motorman. This gave rise to the modern subway train. Following the improvement of an overhead "trolley" system on streetcars for collecting electricity from overhead wires by Sprague, electric tram systems were rapidly adopted across the world.[citation needed]
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+ Earlier electric trains proved difficult or unreliable and experienced limited success until the second half of the 1880s, when new types of current collectors were developed [34]. Siemens' line, for example, provided power through a live rail and a return rail, like a model train, limiting the voltage that could be used, and delivering electric shocks to people and animals crossing the tracks.[38] Siemens later designed his own version of overhead current collection, called the bow collector, and Thorold, Ontario, opened in 1887, and was considered quite successful at the time. While this line proved quite versatile as one of the earliest fully functional electric streetcar installations, it required horse-drawn support while climbing the Niagara Escarpment and for two months of the winter when hydroelectricity was not available. It continued in service in its original form into the 1950s.[citation needed]
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+
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+ Sidney Howe Short designed and produced the first electric motor that operated a streetcar without gears. The motor had its armature direct-connected to the streetcar's axle for the driving force.[39][40][41][42][43] Short pioneered "use of a conduit system of concealed feed" thereby eliminating the necessity of overhead wire and a trolley pole for street cars and railways.[44][39][40] While at the University of Denver he conducted important experiments which established that multiple unit powered cars were a better way to operate trains and trolleys.[39][40]
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+ Electric tramways spread to many European cities in the 1890s, such as Prague, Bohemia (then in the Austro-Hungarian Empire), in 1891; Kiev, Ukraine, in 1892 (the first permanent electric tram line in the Russian Empire); Dresden, Germany, Lyon, France, and Milan and Genoa, Italy, in 1893; Rome, Italy, Plauen, Germany, in 1894; Bristol, United Kingdom, Munich, in 1895; Bilbao, Spain, in 1896; Copenhagen, Denmark, and Vienna, Austria, in 1897; Florence and Turin, Italy, in 1898; Helsinki, Finland, and Madrid and Barcelona, Spain, in 1899.[34] Sarajevo built a citywide system of electric trams in 1895.[45] Budapest established its tramway system in 1887, and its ring line has grown to be the busiest tram line in Europe, with a tram running every 60 seconds at rush hour. Bucharest and Belgrade[46] ran a regular service from 1894.[47][48] Ljubljana introduced its tram system in 1901 – it closed in 1958.[49] Oslo had the first tramway in Scandinavia, starting operation on 2 March 1894.[50]
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+
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+ The first electric tramway in Australia was a Sprague system demonstrated at the 1888 Melbourne Centennial Exhibition in Melbourne; afterwards, this was installed as a commercial venture operating between the outer Melbourne suburb of Box Hill and the then tourist-oriented country town Doncaster from 1889 to 1896.[51] As well, electric systems were built in Adelaide, Ballarat, Bendigo, Brisbane, Fremantle, Geelong, Hobart, Kalgoorlie, Launceston, Leonora, Newcastle, Perth, and Sydney.
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+
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+ By the 1970s, the only full tramway system remaining in Australia was the Melbourne tram system. However, there were also a few single lines remaining elsewhere: the Glenelg tram line, connecting Adelaide to the beachside suburb of Glenelg, and tourist trams in the Victorian Goldfields cities of Bendigo and Ballarat. In recent years the Melbourne system, generally recognised as the largest urban tram network in the world, has been considerably modernised and expanded.[52] The Adelaide line has also been extended to the Entertainment Centre, and work is progressing on further extensions.[53] Sydney re-introduced trams (or light rail) on 31 August 1997. A completely new system, known as G:link, was introduced on the Gold Coast, Queensland on 20 July 2014. The Newcastle Light Rail opened in February 2019, while the Canberra light rail is scheduled to open in April 2019.[54] This will be the first time that there have been trams in Canberra, even though Walter Burley Griffin's 1914-1920 plans for the capital then in the planning stage did propose a Canberra tram system.[55]
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+ In Japan, the Kyoto Electric railroad was the first tram system, starting operation in 1895.[56] By 1932, the network had grown to 82 railway companies in 65 cities, with a total network length of 1,479 km (919 mi).[57] By the 1960s the tram had generally died out in Japan.[58][59]
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+
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+ Two rare but significant alternatives were conduit current collection, which was widely used in London, Washington, D.C. and New York City, and the surface contact collection method, used in Wolverhampton (the Lorain system), Torquay and Hastings in the UK (the Dolter stud system), and currently in Bordeaux, France (the ground-level power supply system).[citation needed]
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+
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+ The convenience and economy of electricity resulted in its rapid adoption once the technical problems of production and transmission of electricity were solved. Electric trams largely replaced animal power and other forms of motive power including cable and steam, in the late 19th and early 20th centuries.[citation needed]
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+ There is one particular hazard associated with trams powered from a trolley pole off an overhead line. Since the tram relies on contact with the rails for the current return path, a problem arises if the tram is derailed or (more usually) if it halts on a section of track that has been particularly heavily sanded by a previous tram, and the tram loses electrical contact with the rails. In this event, the underframe of the tram, by virtue of a circuit path through ancillary loads (such as interior lighting), is live at the full supply voltage, typically 600 volts DC. In British terminology, such a tram was said to be ‘grounded’—not to be confused with the US English use of the term, which means the exact opposite. Any person stepping off the tram completed the earth return circuit and could receive a nasty electric shock. In such an event, the driver was required to jump off the tram (avoiding simultaneous contact with the tram and the ground) and pull down the trolley pole, before allowing passengers off the tram. Unless derailed, the tram could usually be recovered by running water down the running rails from a point higher than the tram, the water providing a conducting bridge between the tram and the rails.[citation needed]
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+ In the 2000s, two companies introduced catenary-free designs. Alstom's Citadis line uses a third rail, and Bombardier's PRIMOVE LRV is charged by contactless induction plates embedded in the trackway.[60]
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+ In some places, other forms of power were used to power the tram.
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+
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+ As early as 1834, Thomas Davenport, a Vermont blacksmith, had invented a battery-powered electric motor which he later patented. The following year he used it to operate a small model electric car on a short section of track four feet in diameter.[61][62]
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+
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+ Attempts to use batteries as a source of electricity were made from the 1880s and 1890s, with unsuccessful trials conducted in among other places Bendigo and Adelaide in Australia, and for about 14 years as The Hague accutram of HTM in the Netherlands. The first trams in Bendigo, Australia, in 1892, were battery-powered but within as little as three months they were replaced with horse-drawn trams. In New York City some minor lines also used storage batteries. Then, comparatively recently, during the 1950s, a longer battery-operated tramway line ran from Milan to Bergamo. In China there is a Nanjing battery Tram line and has been running since 2014.[63] More recently in 2019, the West Midlands Metro in Birmingham, England has adopted battery powered trams on sections through the city centre close to Grade I listed Birmingham Town Hall.
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+ Paris and Berne (Switzerland) [64][circular reference] operated trams that were powered by compressed air using the Mekarski system.
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+ The Convict Tramway [65] was hauled by human power in the form of convicts from the Port Arthur convict settlement.[66] and was created to replace the hazardous sea voyage from Hobart to Port Arthur, Tasmania.[67][65] Charles O'Hara Booth oversaw the construction of the tramway.[68]
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+ It opened in 1836 and ran for 8 km (5 miles) from Oakwood to Taranna.[69] By most definitions, the tramway was the first passenger-carrying railway/tramway in Australia.[67] An unconfirmed report says that it continued to Eaglehawk Neck and, if this was so, the length of the tramway would have been more than doubled. The tramway carried passengers and freight, and ran on wooden rails. The gauge is unknown. The date of closure is unknown, but it was certainly prior to 1877.[70]
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+ In March 2015, China South Rail Corporation (CSR) demonstrated the world's first hydrogen fuel cell vehicle tramcar at an assembly facility in Qingdao. The chief engineer of the CSR subsidiary CSR Sifang Co Ltd., Liang Jianying, said that the company is studying how to reduce the running costs of the tram.[71][72]
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+ The Trieste–Opicina tramway in Trieste operates a hybrid funicular tramway system. Conventional electric trams are operated in street running and on reserved track for most of their route. However, on one steep segment of track, they are assisted by cable tractors, which push the trams uphill and act as brakes for the downhill run. For safety, the cable tractors are always deployed on the downhill side of the tram vehicle.
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+ Similar systems were used elsewhere in the past, notably on the Queen Anne Counterbalance in Seattle and the Darling Street wharf line in Sydney.
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+ Hastings and some other tramways, for example Stockholms Spårvägar in Sweden and some lines in Karachi, used petrol trams. Galveston Island Trolley in Texas operated diesel trams due to the city's hurricane-prone location, which would result in frequent damage to an electrical supply system.
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+ Although Portland, Victoria promotes its tourist tram[73] as being a cable car it actually operates using a hidden diesel motor. The tram, which runs on a circular route around the town of Portland, uses dummies and salons formerly used on the extensive Melbourne cable tramway system and now beautifully restored.
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+ In the mid-20th century many tram systems were disbanded, replaced by buses, automobiles or rapid transit. The General Motors streetcar conspiracy was a case study of the decline of trams in the United States. In the 21st century, trams have been re-introduced in cities where they had been closed down for decades (such as Tramlink in London), or kept in heritage use (such as Spårväg City in Stockholm). Vehicle fabricates from the 1990s and onwards (such as the Bombardier Flexity series and Alstom Citadis) are usually low-floor trams with features such as articulation and regenerative braking.
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+ Trams have been used for two main purposes: for carrying passengers and for carrying cargo. There are several types of passenger tram:
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+ There are two main types of tramways, the classic tramway built in the early 20th century with the tram system operating in mixed traffic, and the later type which is most often associated with the tram system having its own right of way. Tram systems that have their own right of way are often called light rail but this does not always hold true. Though these two systems differ in their operation, their equipment is much the same.
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+ Trams were traditionally operated with separate levers for applying power and brakes. More modern vehicles use a locomotive-style controller which incorporate a dead man's switch. The success of the PCC streetcar had also seen trams use automobile-style foot controls allowing hands-free operation, particularly when the driver was responsible for fare collection.
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+ Electric trams use various devices to collect power from overhead lines. The most common device found today is the pantograph, while some older systems use trolley poles or bow collectors. Ground-level power supply has become a recent innovation. Another new technology uses supercapacitors; when an insulator at a track switch cuts off power from the tram for a short distance along the line, the tram can use energy stored in a large capacitor to drive the tram past the gap in the power feed.[74] A rather obsolete system for power supply is conduit current collection.
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+ The old tram systems in London, Manhattan (New York City), and Washington, D.C., used live rails, like those on third-rail electrified railways, but in a conduit underneath the road, from which they drew power through a plough. It was called Conduit current collection. Washington's was the last of these to close, in 1962. Today, no commercial tramway uses this system. More recently, a modern equivalent to these systems has been developed which allows for the safe installation of a third rail on city streets, which is known as surface current collection or ground-level power supply; the main example of this is the new tramway in Bordeaux.
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+ A ground-level power supply system also known as Surface current collection or Alimentation par le sol (APS) is an updated version of the original stud type system. APS uses a third rail placed between the running rails, divided electrically into eight-metre powered segments with three metre neutral sections between. Each tram has two power collection skates, next to which are antennas that send radio signals to energize the power rail segments as the tram passes over them.
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+ Older systems required mechanical switching systems which were susceptible to environmental problems. At any one time no more than two consecutive segments under the tram should actually be live. Wireless and solid state switching remove the mechanical problem.
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+
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+ Alstom developed the system primarily to avoid intrusive power supply cables in the sensitive area of the old city of old Bordeaux.[75]
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+ Route patterns vary greatly among the world's tram systems, leading to different network topologies.
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+ The resulting route patterns are very different. Some have a rational structure, covering their catchment area as efficiently as possible, with new suburbs being planned with tramlines integral to their layout – such is the case in Amsterdam. Bordeaux and Montpellier have built comprehensive networks, based on radial routes with numerous interconnections, within the last two decades. Some systems serve only parts of their cities, with Berlin being the prime example, owing to the fact that trams survived the city's political division only in the Eastern part. Other systems have ended up with a rather random route map, for instance when some previous operating companies have ceased operation (as with the tramways vicinaux/buurtspoorwegen in Brussels) or where isolated outlying lines have been preserved (as on the eastern fringe of Berlin). In Rome, the remnant of the system comprises 3 isolated radial routes, not connecting in the ancient city centre, but linked by a ring route. Some apparently anomalous lines continue in operation where a new line would not on rational grounds be built, because it is much more costly to build a new line than continue operating an existing one.
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+ In some places, the opportunity is taken when roads are being repaved to lay tramlines (though without erecting overhead cables) even though no service is immediately planned: such is the case in Leipzigerstraße in Berlin, the Haarlemmer Houttuinen in Amsterdam, and Botermarkt in Ghent.
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+ Tram systems operate across national borders in Basel (from Switzerland into France and Germany) and Strasbourg (From France into Germany). It is planned to open a line linking Hasselt (Belgium) with Maastricht (Netherlands) in 2021.
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+ Tramway track can have different rail profiles to accommodate the various operating environments of the vehicle. They may be embedded into concrete for street-running operation, or use standard ballasted track with railroad ties on high-speed sections. A more ecological solution is to embed tracks into grass turf.
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+ Tramway tracks use a grooved rail with a groove designed for tramway or railway track in pavement or grassed surfaces (grassed track or track in a lawn). The rail has the railhead on one side and the guard on the other. The guard provides accommodation for the flange. The guard carries no weight, but may act as a checkrail. Grooved rail was invented in 1852 by Alphonse Loubat, a French inventor who developed improvements in tram and rail equipment, and helped develop tram lines in New York City and Paris. The invention of grooved rail enabled tramways to be laid without causing a nuisance to other road users, except unsuspecting cyclists, who could get their wheels caught in the groove. The grooves may become filled with gravel and dirt (particularly if infrequently used or after a period of idleness) and need clearing from time to time, this being done by a "scrubber" tram. Failure to clear the grooves can lead to a bumpy ride for the passengers, damage to either wheel or rail and possibly derailing.
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+ In narrow situations double-track tram lines sometimes reduce to single track, or, to avoid switches, have the tracks interlaced, e.g. in the Leidsestraat in Amsterdam on three short stretches (see map detail); this is known as interlaced or gauntlet track. There is a UK example of interlaced track on the Tramlink, just west of Mitcham Station, where the formation is narrowed by an old landslip causing an obstruction. (See photo in Tramlink entry).
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+ Historically, the track gauge has had considerable variations, with narrow gauge common in many early systems. However, most light rail systems are now standard gauge. An important advantage of standard gauge is that standard railway maintenance equipment can be used on it, rather than custom-built machinery. Using standard gauge also allows light rail vehicles to be delivered and relocated conveniently using freight railways and locomotives.
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+ Another factor favoring standard gauge is that low-floor vehicles are becoming popular, and there is generally insufficient space for wheelchairs to move between the wheels in a narrow gauge layout. Standard gauge also enables – at least in theory – a larger choice of manufacturers and thus lower procurement costs for new vehicles. However, other factors such as electrification or loading gauge for which there is more variation may require costly custom built units regardless.
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+ Tram stops may be similar to bus stops in design and use, particularly in street-running sections, where in some cases other vehicles are legally required to stop clear of the tram doors. Some stops may resemble to railway platforms, particularly in private right-of-way sections and where trams are boarded at standard railway platform height, as opposed to using steps at the doorway or low-floor trams.
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+ Approximately 5,000 new trams are manufactured each year. As of February 2017, 4,478 new trams were on order from their makers, with options being open for a further 1,092.[76]
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+ The main manufacturers are:
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+ Trams are in a period of growth, with about 800 tram systems operating around the world, 10 or so new systems being opened each year, and many being gradually extended.[87] Some of these systems date from the late 19th or early 20th centuries. In the past 20 years their numbers have been augmented by modern tramway or light rail systems in cities that had discarded this form of transport. There have also been some new tram systems in cities that never previously had them.
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+ Tramways with tramcars (British English) or street railways with streetcars (North American English) were common throughout the industrialised world in the late 19th and early 20th centuries but they had disappeared from most British, Canadian, French and US cities by the mid-20th century.[88]
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+ By contrast, trams in parts of continental Europe continued to be used by many cities, although there were contractions in some countries, including the Netherlands.[89]
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+ Since 1980 trams have returned to favour in many places, partly because their tendency to dominate the roadway, formerly seen as a disadvantage, is now considered to be a merit since it raises the visibility of public transport (encouraging car users to change their mode of travel), and enables streets to be reconfigured to give more space to pedestrians, making cites more pleasant places to live. New systems have been built in the United States, United Kingdom, Ireland, Italy, France, Australia and many other countries.
168
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+ In Milan, Italy, the old "Ventotto" trams are considered by its inhabitants a "symbol" of the city. The same can be said of trams in Melbourne in general, but particularly the iconic W class. The Toronto streetcar system had similarly become an iconic symbol of the city, operating the largest network in the Americas as well as the only large-scale tram system in Canada (not including light rail systems, or heritage lines).[90][91]
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+
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+ The largest tram ((classic tram, streetcar, straßenbahn) and fast tram (light rail, stadtbahn)) networks in the world by route length (as of 2016)[92] are:
172
+
173
+ Other large transit networks that operate streetcar and light rail systems include:
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+
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+ This list is not exhaustive.
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+
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+
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+ Historically, the Paris Tram System was, at its peak, the world's largest system, with 1,111 km (690 mi) of track in 1925[citation needed] (according to other sources, ca. 640 km (400 mi) of route length in 1930). However it was completely closed in 1938.[133] The next largest system appears to have been 857 km (533 mi), in Buenos Aires before 19 February 1963. The third largest was Chicago, with over 850 km (530 mi) of track,[134] but it was all converted to trolleybus and bus services by 21 June 1958. Before its decline, the BVG in Berlin operated a very large network with 634 km (394 mi) of route. Before its system started to be converted to trolleybus (and later bus) services in the 1930s (last tramway closed 6 July 1952), the first-generation London network had 555 km (345 mi) of route in 1931.[135] In 1958 trams in Rio de Jainero were employed on (433 km; 269 mi) of track. The final line, the Santa teresa route was closed in 1968.[136] During a period in the 1980s, the world's largest tram system was in Leningrad (now known as St. Petersburg) with 350 km (220 mi), USSR, and was included as such in the Guinness World Records;[citation needed] however Saint Petersburg's tram system has declined in size since the fall of the Soviet Union. Vienna in 1960 had 340 km (211 mi), before the expansion of bus services and the opening of a subway (1976). Substituting subway services for tram routes continues. 320 km (199 mi) was in Minneapolis-Saint Paul in 1947: There streetcars ended 31 October 1953 in Minneapolis and 19 June 1954 in St. Paul.[137] The Sydney tram network, before it was closed on 25 February 1961, had 291 km (181 mi) of route, and was thus the largest in Australia. As from 1961, the Melbourne system (currently recognised as the world's largest) took over Sydney's title as the largest network in Australia.
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+ In many European cities, much tramway infrastructure was lost in the mid-20th century, though not always on the same scale as in other parts of the world such as North America. Most of Central and Eastern Europe retained the majority of its tramway systems and it is here that the largest and busiest tram systems in the world are found.
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+ Whereas most systems and vehicles in the tram sector are found in Central and Eastern Europe, in the 1960s and 1970s, tram systems were shut down in many places in Western Europe, however urban transportation has been experiencing a sustained long running revival since the 1990s. Many European cities are rehabilitating, upgrading, expanding and reconstructing their old tramway lines and building new tramway lines.[139]
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+ In North America, these vehicles are called "streetcars" (or "trolleys"); the term tram is more likely to be understood as an aerial tramway or a people-mover. Streetcar systems were developed in late 19th to early 20th centuries in a number of cities throughout North America. However, most North American cities saw its streetcar lines removed in the mid-20th century for a variety of financial, technological and social reasons. Exceptions included Boston,[140] Cleveland, Mexico City, New Orleans, Newark, Philadelphia, Pittsburgh, San Francisco, and Toronto.
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+ Toronto currently operates the largest streetcar system in the Americas in terms of track length and ridership. Operated by the Toronto Transit Commission, the streetcar system is the only large-scale streetcar system existing in Canada, excluding heritage streetcar, or light rail systems that are operated in other Canadian municipalities. The streetcar system was established in 1861, and used a variety of vehicles in its history, including horse-drawn streetcars, Peter Witt streetcars, the PCC streetcar, and the Canadian Light Rail Vehicle and its articulated counterpart, the Articulated Light Rail Vehicle. Since December 29, 2019,[141] the system exclusively uses the Flexity Outlook made by Bombardier Transportation.[142][143][144][145]
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+ Streetcars once existed in the Canadian cities of Calgary, Edmonton, Halifax, Hamilton, Kingston, Kitchener, London, Montreal, Ottawa, Peterborough, Quebec City, Regina, Saskatoon, Windsor, and Vancouver. However, Canadian cities excluding Toronto, removed their streetcar systems in the mid-20th century. In the late 1970s and early 1980s, light rail systems were introduced in Calgary and Edmonton; with another light rail system established in Ottawa in 2001. There is now something of a renaissance for light railways in mid-sized cities with Waterloo, Ontario the first to come on line and construction underway in Mississauga, Ontario. In the late 20th century, several Canadian locales restored portions of their defunct streetcar lines, operating them as a heritage feature for tourists. Heritage streetcar lines in Canada include the High Level Bridge Streetcar in Edmonton, the Nelson Electric Tramway in Nelson, and the Whitehorse Waterfront Trolley in Whitehorse.
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+ Pittsburgh had kept most of its streetcar system serving the city and many suburbs, making it the longest-lasting large-network streetcar system in the United States.[citation needed] However, most of Pittsburgh's surviving streetcar lines were converted to light rail in the 1980s. San Francisco's Muni Metro system is the largest surviving streetcar system in the United States, and has even revived previously closed streetcar lines such as the F Market & Wharves heritage streetcar line. In the late 20th century, several cities installed modern light rail systems, in part along the same corridors as their old streetcars systems, the first of these being the San Diego Trolley in San Diego in 1981.
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+ In the 1980s, some cities in the United States brought back streetcars lines, including Memphis, Tampa, and Little Rock; However, these streetcar systems were designed as heritage streetcar lines, and used vintage or replica-vintage vehicles. The first "second-generation streetcar systems" in North America was opened in Portland in 2001.[146] The "second-generation streetcar system," utilizes modern vehicles – vehicles that feature low-floor streetcars. These newer streetcar systems were built in several American cities in the early 21st century including Atlanta, Charlotte, Cincinnati, Dallas, Detroit, Kansas City, Milwaukee, Oklahoma City, Seattle, Tucson, and Washington, D.C..
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+ Model trams are popular in HO scale (1:87) and O scale (1:48 in the US and generally 1:43,5 and 1:45 in Europe and Asia). They are typically powered and will accept plastic figures inside. Common manufacturers are Roco and Lima, with many custom models being made as well. The German firm Hödl[165] and the Austrian Halling[166] specialise in 1:87 scale.[167]
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+ In the US, Bachmann Industries is a mass supplier of HO streetcars and kits. Bowser Manufacturing has produced white metal models for over 50 years.[168] There are many boutique vendors offering limited run epoxy and wood models. At the high end are highly detailed brass models which are usually imported from Japan or Korea and can cost in excess of $500. Many of these run on 16.5 mm (0.65 in) gauge track, which is correct for the representation of 4 ft 8 1⁄2 in (1,435 mm) (standard gauge) in HO scale as in US and Japan, but incorrect in 4 mm (1:76.2) scale, as it represents 4 ft 8 1⁄2 in (1,435 mm). This scale/gauge hybrid is called OO scale.
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+ O scale trams are also very popular among tram modellers because the increased size allows for more detail and easier crafting of overhead wiring. In the US these models are usually purchased in epoxy or wood kits and some as brass models. The Saint Petersburg Tram Company[169] produces highly detailed polyurethane non-powered O Scale models from around the world which can easily be powered by trucks from vendors like Q-Car.[170]
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+ In the US, one of the best resources for model tram enthusiasts is the East Penn Traction Club of Philadelphia [171] and Trolleyville a website of the Southern California Traction Club.[172]
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+ It is thought that the first example of a working model tramcar in the UK built by an amateur for fun was in 1929, when Frank E. Wilson created a replica of London County Council Tramways E class car 444 in 1:16 scale, which he demonstrated at an early Model Engineer Exhibition. Another of his models was London E/1 1800, which was the only tramway exhibit in the Faraday Memorial Exhibition of 1931. Together with likeminded friends, Frank Wilson went on to found the Tramway & Light Railway Society[173] in 1938, establishing tramway modelling as a hobby.
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+ Mr. Nathan was a passenger by No. 2 tramway car [...] [he] alighted from the car at the southern end, but before he got clear of the rails the car moved onwards [...] he was thus whirled round by the sudden motion of the carriage and his body was brought under the front wheel.
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+ A tram (in North America streetcar or trolley) is a rail vehicle that runs on tramway tracks along public urban streets; some include segments of segregated right-of-way.[1][2] The lines or networks operated by tramcars are called tramways. Historically the term electric street railways was also used in the United States. In the United States, the term tram has sometimes been used for rubber-tired trackless trains, which are unrelated to other kinds of trams.
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+ Tram vehicles are usually lighter and shorter than main line and rapid transit trains. Today, most trams use electrical power, usually fed by a pantograph sliding on an overhead line; older systems may use a trolley pole or a bow collector. In some cases, a contact shoe on a third rail is used. If necessary, they may have dual power systems—electricity in city streets and diesel in more rural environments. Occasionally, trams also carry freight.
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+ Trams are now commonly included in the wider term "light rail",[3] which also includes grade-separated systems. Some trams, known as tram-trains, may have segments that run on mainline railway tracks, similar to interurban systems. The differences between these modes of rail transport are often indistinct and a given system may combine multiple features.
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+ One of the advantages over earlier forms of transit was the low rolling resistance of metal wheels on steel rails, allowing the trams to haul a greater load for a given effort. Problems included the high total cost of ownership of horses. Electric trams largely replaced animal power in the late 19th and early 20th centuries. Improvements in other vehicles such as buses led to decline of trams in the mid 20th century. However, trams have seen resurgence in recent years.
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+
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+ The English terms tram and tramway are derived from the Scots word tram,[4] referring respectively to a type of truck (goods wagon or freight railroad car) used in coal mines and the tracks on which they ran. The word tram probably derived from Middle Flemish trame ("beam, handle of a barrow, bar, rung"). The identical word la trame with the meaning "crossbeam" is also used in the French language. Etymologists believe that the word tram refers to the wooden beams the railway tracks were initially made of before the railroad pioneers switched to the much more wear-resistant tracks made of iron and, later, steel.[5] The word Tram-car is attested from 1873.[6]
12
+
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+ Although the terms tram and tramway have been adopted by many languages, they are not used universally in English; North Americans prefer streetcar, trolley, or trolleycar. The term streetcar is first recorded in 1840, and originally referred to horsecars. When electrification came, Americans began to speak of trolleycars or later, trolleys. A widely held belief holds the word to derive from the troller (said to derive from the words traveler and roller), a four-wheeled device that was dragged along dual overhead wires by a cable that connected the troller to the top of the car and collected electrical power from the overhead wires;[7] this portmanteau derivation is, however, most likely folk etymology. "Trolley" and variants refer to the verb troll, meaning "roll" and probably derived from Old French,[8] and cognate uses of the word were well established for handcarts and horse drayage, as well as for nautical uses.[9]
14
+
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+ The alternative North American term 'trolley' may strictly speaking be considered incorrect, as the term can also be applied to cable cars, or conduit cars that instead draw power from an underground supply. Conventional diesel tourist buses decorated to look like streetcars are sometimes called trolleys in the US (tourist trolley). Furthering confusion, the term tram has instead been applied to open-sided, low-speed segmented vehicles on rubber tires generally used to ferry tourists short distances, for example on the Universal Studios backlot tour and, in many countries, as tourist transport to major destinations. The term may also apply to an aerial ropeway, e.g. the Roosevelt Island Tramway.
16
+
17
+ Although the use of the term trolley for tram was not adopted in Europe, the term was later associated with the trolleybus, a rubber-tyred vehicle running on hard pavement, which draws its power from pairs of overhead wires. These electric buses, which use twin trolley poles, are also called trackless trolleys (particularly in the northeastern US), or sometimes simply trolleys (in the UK, as well as the Pacific Northwest, including Seattle, and Vancouver).
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+
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+ The New South Wales government in Australia has decided to use the term "light rail" for their trams.
20
+
21
+ The history of trams, streetcars or trolley systems, began in the early nineteenth century. It can be divided up into several discrete periods defined by the principal means of motive power used.
22
+
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+ The world's first passenger train or tram was the Swansea and Mumbles Railway, in Wales, UK. The Mumbles Railway Act was passed by the British Parliament in 1804, and horse-drawn service started in 1807.[10] The service closed in 1827, but was restarted in 1860, again using horses.[11] It was worked by steam from 1877, and then, from 1929, by very large (106-seater) electric tramcars, until closure in 1961.[citation needed] The Swansea and Mumbles Railway was something of a one-off however, and no street tramway would appear in Britain until 1860 when one was built in Birkenhead by the American George Francis Train.[12]
24
+
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+ Street railways developed in America before Europe, largely due to the poor paving of the streets in American cities which made them unsuitable for horsebuses, which were then common on the well-paved streets of European cities. Running the horsecars on rails allowed for a much smoother ride. There are records of a street railway running in Baltimore as early as 1828, however the first authenticated streetcar in America, was the New York and Harlem Railroad developed by the Irish coach builder John Stephenson, in New York City which began service in the year 1832.[13][14] The New York and Harlem Railroad's Fourth Avenue Line ran along the Bowery and Fourth Avenue in New York City. It was followed in 1835 by the New Orleans and Carrollton Railroad in New Orleans, Louisiana,[15] which still operates as the St. Charles Streetcar Line. Other American cities did not follow until the 1850s, after which the "animal railway" became an increasingly common feature in the larger towns.[15]
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+
27
+ The first permanent tram line in continental Europe was opened in Paris in 1855 by Alphonse Loubat who had previously worked on American streetcar lines.[16] The tram was developed in numerous cities of Europe (some of the most extensive systems were found in Berlin, Budapest, Birmingham, Leningrad, Lisbon, London, Manchester, Paris).
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+
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+ The first tram in South America opened in 1858 in Santiago, Chile. The first trams in Australia opened in 1860 in Sydney. Africa's first tram service started in Alexandria on 8 January 1863. The first trams in Asia opened in 1869 in Batavia (now Jakarta), Netherlands East Indies (now Indonesia).
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+
31
+ Problems with horsecars included the fact that any given animal could only work so many hours on a given day, had to be housed, groomed, fed and cared for day in and day out, and produced prodigious amounts of manure, which the streetcar company was charged with storing and then disposing of. Since a typical horse pulled a streetcar for about a dozen miles a day and worked for four or five hours, many systems needed ten or more horses in stable for each horsecar.
32
+
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+ Horsecars were largely replaced by electric-powered trams following the improvement of an overhead trolley system on trams for collecting electricity from overhead wires by Frank J. Sprague. His spring-loaded trolley pole used a wheel to travel along the wire. In late 1887 and early 1888, using his trolley system, Sprague installed the first successful large electric street railway system in Richmond, Virginia. Within a year, the economy of electric power had replaced more costly horsecars in many cities. By 1889, 110 electric railways incorporating Sprague's equipment had been begun or planned on several continents.[17]
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+
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+ Horses continued to be used for light shunting well into the 20th century, and many large metropolitan lines lasted into the early 20th century. New York City had a regular horsecar service on the Bleecker Street Line until its closure in 1917.[18] Pittsburgh, Pennsylvania, had its Sarah Street line drawn by horses until 1923. The last regular mule-drawn cars in the US ran in Sulphur Rock, Arkansas, until 1926 and were commemorated by a U.S. postage stamp issued in 1983.[19] The last mule tram service in Mexico City ended in 1932, and a mule tram in Celaya, Mexico, survived until 1954.[20] The last horse-drawn tram to be withdrawn from public service in the UK took passengers from Fintona railway station to Fintona Junction one mile away on the main Omagh to Enniskillen railway in Northern Ireland. The tram made its last journey on 30 September 1957 when the Omagh to Enniskillen line closed. The "van" now lies at the Ulster Transport Museum.
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+ Horse-drawn trams still operate on the 1876-built Douglas Bay Horse Tramway in the Isle of Man, and at the 1894-built horse tram at Victor Harbor in South Australia. New horse-drawn systems have been established at the Hokkaidō Museum in Japan and also in Disneyland. A horse tram route in Polish gmina Mrozy, first built in 1902, was reopened in 2012.
38
+
39
+ The first mechanical trams were powered by steam. Generally, there were two types of steam tram. The first and most common had a small steam locomotive (called a tram engine in the UK) at the head of a line of one or more carriages, similar to a small train. Systems with such steam trams included Christchurch, New Zealand; Sydney, Australia; other city systems in New South Wales; Munich, Germany (from August 1883 on),[21] British India (Pakistan) (from 1885) and the Dublin & Blessington Steam Tramway (from 1888) in Ireland. Steam tramways also were used on the suburban tramway lines around Milan and Padua; the last Gamba de Legn ("Peg-Leg") tramway ran on the Milan-Magenta-Castano Primo route in late 1957.[22]
40
+
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+ The other style of steam tram had the steam engine in the body of the tram, referred to as a tram engine (UK) or steam dummy (US). The most notable system to adopt such trams was in Paris. French-designed steam trams also operated in Rockhampton, in the Australian state of Queensland between 1909 and 1939. Stockholm, Sweden, had a steam tram line at the island of Södermalm between 1887 and 1901.
42
+
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+ Tram engines usually had modifications to make them suitable for street running in residential areas. The wheels, and other moving parts of the machinery, were usually enclosed for safety reasons and to make the engines quieter. Measures were often taken to prevent the engines from emitting visible smoke or steam. Usually the engines used coke rather than coal as fuel to avoid emitting smoke; condensers or superheating were used to avoid emitting visible steam. A major drawback of this style of tram was the limited space for the engine, so that these trams were usually underpowered. Steam tram engines faded out around the 1890s to 1900s, being replaced by electric trams.
44
+
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+ Another motive system for trams was the cable car, which was pulled along a fixed track by a moving steel cable. The power to move the cable was normally provided at a "powerhouse" site a distance away from the actual vehicle. The London and Blackwall Railway, which opened for passengers in east London, England, in 1840 used such a system.[23]
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+
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+ The first practical cable car line was tested in San Francisco, in 1873. Part of its success is attributed to the development of an effective and reliable cable grip mechanism, to grab and release the moving cable without damage. The second city to operate cable trams was Dunedin in New Zealand, from 1881 to 1957.
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+ The most extensive cable system in the US was built in Chicago, having been built in stages between 1859 and 1892. New York City developed multiple cable car lines, that operated from 1883 to 1909.[24] Los Angeles also had several cable car lines, including the Second Street Cable Railroad, which operated from 1885 to 1889, and the Temple Street Cable Railway, which operated from 1886 to 1898.
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+
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+ From 1885 to 1940, the city of Melbourne, Victoria, Australia operated one of the largest cable systems in the world, at its peak running 592 trams on 75 kilometres (47 mi) of track. There were also two isolated cable lines in Sydney, New South Wales, Australia; the North Sydney line from 1886 to 1900,[25] and the King Street line from 1892 to 1905.
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+
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+ In Dresden, Germany, in 1901 an elevated suspended cable car following the Eugen Langen one-railed floating tram system started operating. Cable cars operated on Highgate Hill in North London and Kennington to Brixton Hill In South London.[when?] They also worked around "Upper Douglas" in the Isle of Man from 1897 to 1929 (cable car 72/73 is the sole survivor of the fleet).
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+
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+ Cable cars suffered from high infrastructure costs, since an expensive system of cables, pulleys, stationary engines and lengthy underground vault structures beneath the rails had to be provided. They also required physical strength and skill to operate, and alert operators to avoid obstructions and other cable cars. The cable had to be disconnected ("dropped") at designated locations to allow the cars to coast by inertia, for example when crossing another cable line. The cable would then have to be "picked up" to resume progress, the whole operation requiring precise timing to avoid damage to the cable and the grip mechanism. Breaks and frays in the cable, which occurred frequently, required the complete cessation of services over a cable route while the cable was repaired. Due to overall wear, the entire length of cable (typically several kilometres) would have to be replaced on a regular schedule. After the development of reliable electrically powered trams, the costly high-maintenance cable car systems were rapidly replaced in most locations.
56
+
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+ Cable cars remained especially effective in hilly cities, since their nondriven wheels would not lose traction as they climbed or descended a steep hill. The moving cable would physically pull the car up the hill at a steady pace, unlike a low-powered steam or horse-drawn car. Cable cars do have wheel brakes and track brakes, but the cable also helps restrain the car to going downhill at a constant speed. Performance in steep terrain partially explains the survival of cable cars in San Francisco.
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+
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+ The San Francisco cable cars, though significantly reduced in number, continue to perform a regular transportation function, in addition to being a well-known tourist attraction. A single cable line also survives in Wellington, New Zealand (rebuilt in 1979 as a funicular but still called the "Wellington Cable Car"). Another system, actually two separate cable lines with a shared power station in the middle, operates from the Welsh town of Llandudno up to the top of the Great Orme hill in North Wales, UK.
60
+
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+ In the late 19th and early 20th centuries a number of systems in various parts of the world employed trams powered by gas, naphtha gas or coal gas in particular. Gas trams are known to have operated between Alphington and Clifton Hill in the northern suburbs of Melbourne, Australia (1886–1888); in Berlin and Dresden, Germany; in Estonia (1921–1951); between Jelenia Góra, Cieplice, and Sobieszów in Poland (from 1897); and in the UK at Lytham St Annes, Neath (1896–1920), and Trafford Park, Manchester (1897–1908).
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+
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+ On 29 December 1886 the Melbourne newspaper The Argus reprinted a report from the San Francisco Bulletin that Mr Noble had demonstrated a new 'motor car' for tramways 'with success'. The tramcar 'exactly similar in size, shape, and capacity to a cable grip car' had the 'motive power' of gas 'with which the reservoir is to be charged once a day at power stations by means of a rubber hose'. The car also carried an electricity generator for 'lighting up the tram and also for driving the engine on steep grades and effecting a start'.[26]
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+ Comparatively little has been published about gas trams. However, research on the subject was carried out for an article in the October 2011 edition of "The Times", the historical journal of the Australian Association of Timetable Collectors, now the Australian Timetable Association.[27][28][29][30]
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+ A tram system powered by compressed natural gas was due to open in Malaysia in 2012,[31] but the news about the project appears to have dried up.
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+ The world's first electric tram line operated in Sestroretsk near Saint Petersburg, Russia, invented and tested by Fyodor Pirotsky in 1875.[32][33] Later, using a similar technology, Pirotsky put into service the first public electric tramway in St. Petersburg, which operated only during September 1880.[34]
70
+ The second demonstrative tramway was presented by Siemens & Halske at the 1879 Berlin Industrial Exposition.
71
+ The first public electric tramway used for permanent service was the Gross-Lichterfelde tramway in Lichterfelde near Berlin in Germany, which opened in 1881. It was built by Werner von Siemens who contacted Pirotsky. This was world's first commercially successful electric tram. It initially drew current from the rails, with overhead wire being installed in 1883.[35]
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+ In Britain, Volk's Electric Railway was opened in 1883 in Brighton). This two kilometer line along the seafront, re-gauged to 2 ft 9 in (838 mm) in 1884, remains in service to this day and is the oldest operating electric tramway in the world. Also in 1883, Mödling and Hinterbrühl Tram was opened near Vienna in Austria. It was the first tram in the world in regular service that was run with electricity served by an overhead line with pantograph current collectors. The Blackpool Tramway was opened in Blackpool, UK on 29 September 1885 using conduit collection along Blackpool Promenade. This system is still in operation in a modernised form.[36]
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+
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+ Earliest tram system in Canada was by John Joseph Wright, brother of the famous mining entrepreneur Whitaker Wright, in Toronto in 1883, introducing electric trams in 1892. In the US, multiple functioning experimental electric trams were exhibited at the 1884 World Cotton Centennial World's Fair in New Orleans, Louisiana, but they were not deemed good enough to replace the Lamm fireless engines then propelling the St. Charles Avenue Streetcar in that city. The first commercial installation of an electric streetcar in the United States was built in 1884 in Cleveland, Ohio and operated for a period of one year by the East Cleveland Street Railway Company.[37] Trams were operated in Richmond, Virginia, in 1888, on the Richmond Union Passenger Railway built by Frank J. Sprague. Sprague later developed multiple unit control, first demonstrated in Chicago in 1897, allowing multiple cars to be coupled together and operated by a single motorman. This gave rise to the modern subway train. Following the improvement of an overhead "trolley" system on streetcars for collecting electricity from overhead wires by Sprague, electric tram systems were rapidly adopted across the world.[citation needed]
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+
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+ Earlier electric trains proved difficult or unreliable and experienced limited success until the second half of the 1880s, when new types of current collectors were developed [34]. Siemens' line, for example, provided power through a live rail and a return rail, like a model train, limiting the voltage that could be used, and delivering electric shocks to people and animals crossing the tracks.[38] Siemens later designed his own version of overhead current collection, called the bow collector, and Thorold, Ontario, opened in 1887, and was considered quite successful at the time. While this line proved quite versatile as one of the earliest fully functional electric streetcar installations, it required horse-drawn support while climbing the Niagara Escarpment and for two months of the winter when hydroelectricity was not available. It continued in service in its original form into the 1950s.[citation needed]
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+
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+ Sidney Howe Short designed and produced the first electric motor that operated a streetcar without gears. The motor had its armature direct-connected to the streetcar's axle for the driving force.[39][40][41][42][43] Short pioneered "use of a conduit system of concealed feed" thereby eliminating the necessity of overhead wire and a trolley pole for street cars and railways.[44][39][40] While at the University of Denver he conducted important experiments which established that multiple unit powered cars were a better way to operate trains and trolleys.[39][40]
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+ Electric tramways spread to many European cities in the 1890s, such as Prague, Bohemia (then in the Austro-Hungarian Empire), in 1891; Kiev, Ukraine, in 1892 (the first permanent electric tram line in the Russian Empire); Dresden, Germany, Lyon, France, and Milan and Genoa, Italy, in 1893; Rome, Italy, Plauen, Germany, in 1894; Bristol, United Kingdom, Munich, in 1895; Bilbao, Spain, in 1896; Copenhagen, Denmark, and Vienna, Austria, in 1897; Florence and Turin, Italy, in 1898; Helsinki, Finland, and Madrid and Barcelona, Spain, in 1899.[34] Sarajevo built a citywide system of electric trams in 1895.[45] Budapest established its tramway system in 1887, and its ring line has grown to be the busiest tram line in Europe, with a tram running every 60 seconds at rush hour. Bucharest and Belgrade[46] ran a regular service from 1894.[47][48] Ljubljana introduced its tram system in 1901 – it closed in 1958.[49] Oslo had the first tramway in Scandinavia, starting operation on 2 March 1894.[50]
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+
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+ The first electric tramway in Australia was a Sprague system demonstrated at the 1888 Melbourne Centennial Exhibition in Melbourne; afterwards, this was installed as a commercial venture operating between the outer Melbourne suburb of Box Hill and the then tourist-oriented country town Doncaster from 1889 to 1896.[51] As well, electric systems were built in Adelaide, Ballarat, Bendigo, Brisbane, Fremantle, Geelong, Hobart, Kalgoorlie, Launceston, Leonora, Newcastle, Perth, and Sydney.
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+
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+ By the 1970s, the only full tramway system remaining in Australia was the Melbourne tram system. However, there were also a few single lines remaining elsewhere: the Glenelg tram line, connecting Adelaide to the beachside suburb of Glenelg, and tourist trams in the Victorian Goldfields cities of Bendigo and Ballarat. In recent years the Melbourne system, generally recognised as the largest urban tram network in the world, has been considerably modernised and expanded.[52] The Adelaide line has also been extended to the Entertainment Centre, and work is progressing on further extensions.[53] Sydney re-introduced trams (or light rail) on 31 August 1997. A completely new system, known as G:link, was introduced on the Gold Coast, Queensland on 20 July 2014. The Newcastle Light Rail opened in February 2019, while the Canberra light rail is scheduled to open in April 2019.[54] This will be the first time that there have been trams in Canberra, even though Walter Burley Griffin's 1914-1920 plans for the capital then in the planning stage did propose a Canberra tram system.[55]
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+
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+ In Japan, the Kyoto Electric railroad was the first tram system, starting operation in 1895.[56] By 1932, the network had grown to 82 railway companies in 65 cities, with a total network length of 1,479 km (919 mi).[57] By the 1960s the tram had generally died out in Japan.[58][59]
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+ Two rare but significant alternatives were conduit current collection, which was widely used in London, Washington, D.C. and New York City, and the surface contact collection method, used in Wolverhampton (the Lorain system), Torquay and Hastings in the UK (the Dolter stud system), and currently in Bordeaux, France (the ground-level power supply system).[citation needed]
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+
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+ The convenience and economy of electricity resulted in its rapid adoption once the technical problems of production and transmission of electricity were solved. Electric trams largely replaced animal power and other forms of motive power including cable and steam, in the late 19th and early 20th centuries.[citation needed]
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+
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+ There is one particular hazard associated with trams powered from a trolley pole off an overhead line. Since the tram relies on contact with the rails for the current return path, a problem arises if the tram is derailed or (more usually) if it halts on a section of track that has been particularly heavily sanded by a previous tram, and the tram loses electrical contact with the rails. In this event, the underframe of the tram, by virtue of a circuit path through ancillary loads (such as interior lighting), is live at the full supply voltage, typically 600 volts DC. In British terminology, such a tram was said to be ‘grounded’—not to be confused with the US English use of the term, which means the exact opposite. Any person stepping off the tram completed the earth return circuit and could receive a nasty electric shock. In such an event, the driver was required to jump off the tram (avoiding simultaneous contact with the tram and the ground) and pull down the trolley pole, before allowing passengers off the tram. Unless derailed, the tram could usually be recovered by running water down the running rails from a point higher than the tram, the water providing a conducting bridge between the tram and the rails.[citation needed]
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+
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+ In the 2000s, two companies introduced catenary-free designs. Alstom's Citadis line uses a third rail, and Bombardier's PRIMOVE LRV is charged by contactless induction plates embedded in the trackway.[60]
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+
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+ In some places, other forms of power were used to power the tram.
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+
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+ As early as 1834, Thomas Davenport, a Vermont blacksmith, had invented a battery-powered electric motor which he later patented. The following year he used it to operate a small model electric car on a short section of track four feet in diameter.[61][62]
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+
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+ Attempts to use batteries as a source of electricity were made from the 1880s and 1890s, with unsuccessful trials conducted in among other places Bendigo and Adelaide in Australia, and for about 14 years as The Hague accutram of HTM in the Netherlands. The first trams in Bendigo, Australia, in 1892, were battery-powered but within as little as three months they were replaced with horse-drawn trams. In New York City some minor lines also used storage batteries. Then, comparatively recently, during the 1950s, a longer battery-operated tramway line ran from Milan to Bergamo. In China there is a Nanjing battery Tram line and has been running since 2014.[63] More recently in 2019, the West Midlands Metro in Birmingham, England has adopted battery powered trams on sections through the city centre close to Grade I listed Birmingham Town Hall.
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+
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+ Paris and Berne (Switzerland) [64][circular reference] operated trams that were powered by compressed air using the Mekarski system.
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+
105
+ The Convict Tramway [65] was hauled by human power in the form of convicts from the Port Arthur convict settlement.[66] and was created to replace the hazardous sea voyage from Hobart to Port Arthur, Tasmania.[67][65] Charles O'Hara Booth oversaw the construction of the tramway.[68]
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+
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+ It opened in 1836 and ran for 8 km (5 miles) from Oakwood to Taranna.[69] By most definitions, the tramway was the first passenger-carrying railway/tramway in Australia.[67] An unconfirmed report says that it continued to Eaglehawk Neck and, if this was so, the length of the tramway would have been more than doubled. The tramway carried passengers and freight, and ran on wooden rails. The gauge is unknown. The date of closure is unknown, but it was certainly prior to 1877.[70]
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+
109
+ In March 2015, China South Rail Corporation (CSR) demonstrated the world's first hydrogen fuel cell vehicle tramcar at an assembly facility in Qingdao. The chief engineer of the CSR subsidiary CSR Sifang Co Ltd., Liang Jianying, said that the company is studying how to reduce the running costs of the tram.[71][72]
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+
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+ The Trieste–Opicina tramway in Trieste operates a hybrid funicular tramway system. Conventional electric trams are operated in street running and on reserved track for most of their route. However, on one steep segment of track, they are assisted by cable tractors, which push the trams uphill and act as brakes for the downhill run. For safety, the cable tractors are always deployed on the downhill side of the tram vehicle.
112
+
113
+ Similar systems were used elsewhere in the past, notably on the Queen Anne Counterbalance in Seattle and the Darling Street wharf line in Sydney.
114
+
115
+ Hastings and some other tramways, for example Stockholms Spårvägar in Sweden and some lines in Karachi, used petrol trams. Galveston Island Trolley in Texas operated diesel trams due to the city's hurricane-prone location, which would result in frequent damage to an electrical supply system.
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+
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+ Although Portland, Victoria promotes its tourist tram[73] as being a cable car it actually operates using a hidden diesel motor. The tram, which runs on a circular route around the town of Portland, uses dummies and salons formerly used on the extensive Melbourne cable tramway system and now beautifully restored.
118
+
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+ In the mid-20th century many tram systems were disbanded, replaced by buses, automobiles or rapid transit. The General Motors streetcar conspiracy was a case study of the decline of trams in the United States. In the 21st century, trams have been re-introduced in cities where they had been closed down for decades (such as Tramlink in London), or kept in heritage use (such as Spårväg City in Stockholm). Vehicle fabricates from the 1990s and onwards (such as the Bombardier Flexity series and Alstom Citadis) are usually low-floor trams with features such as articulation and regenerative braking.
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+
121
+ Trams have been used for two main purposes: for carrying passengers and for carrying cargo. There are several types of passenger tram:
122
+
123
+ There are two main types of tramways, the classic tramway built in the early 20th century with the tram system operating in mixed traffic, and the later type which is most often associated with the tram system having its own right of way. Tram systems that have their own right of way are often called light rail but this does not always hold true. Though these two systems differ in their operation, their equipment is much the same.
124
+
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+ Trams were traditionally operated with separate levers for applying power and brakes. More modern vehicles use a locomotive-style controller which incorporate a dead man's switch. The success of the PCC streetcar had also seen trams use automobile-style foot controls allowing hands-free operation, particularly when the driver was responsible for fare collection.
126
+
127
+ Electric trams use various devices to collect power from overhead lines. The most common device found today is the pantograph, while some older systems use trolley poles or bow collectors. Ground-level power supply has become a recent innovation. Another new technology uses supercapacitors; when an insulator at a track switch cuts off power from the tram for a short distance along the line, the tram can use energy stored in a large capacitor to drive the tram past the gap in the power feed.[74] A rather obsolete system for power supply is conduit current collection.
128
+
129
+ The old tram systems in London, Manhattan (New York City), and Washington, D.C., used live rails, like those on third-rail electrified railways, but in a conduit underneath the road, from which they drew power through a plough. It was called Conduit current collection. Washington's was the last of these to close, in 1962. Today, no commercial tramway uses this system. More recently, a modern equivalent to these systems has been developed which allows for the safe installation of a third rail on city streets, which is known as surface current collection or ground-level power supply; the main example of this is the new tramway in Bordeaux.
130
+
131
+ A ground-level power supply system also known as Surface current collection or Alimentation par le sol (APS) is an updated version of the original stud type system. APS uses a third rail placed between the running rails, divided electrically into eight-metre powered segments with three metre neutral sections between. Each tram has two power collection skates, next to which are antennas that send radio signals to energize the power rail segments as the tram passes over them.
132
+
133
+ Older systems required mechanical switching systems which were susceptible to environmental problems. At any one time no more than two consecutive segments under the tram should actually be live. Wireless and solid state switching remove the mechanical problem.
134
+
135
+ Alstom developed the system primarily to avoid intrusive power supply cables in the sensitive area of the old city of old Bordeaux.[75]
136
+
137
+ Route patterns vary greatly among the world's tram systems, leading to different network topologies.
138
+
139
+ The resulting route patterns are very different. Some have a rational structure, covering their catchment area as efficiently as possible, with new suburbs being planned with tramlines integral to their layout – such is the case in Amsterdam. Bordeaux and Montpellier have built comprehensive networks, based on radial routes with numerous interconnections, within the last two decades. Some systems serve only parts of their cities, with Berlin being the prime example, owing to the fact that trams survived the city's political division only in the Eastern part. Other systems have ended up with a rather random route map, for instance when some previous operating companies have ceased operation (as with the tramways vicinaux/buurtspoorwegen in Brussels) or where isolated outlying lines have been preserved (as on the eastern fringe of Berlin). In Rome, the remnant of the system comprises 3 isolated radial routes, not connecting in the ancient city centre, but linked by a ring route. Some apparently anomalous lines continue in operation where a new line would not on rational grounds be built, because it is much more costly to build a new line than continue operating an existing one.
140
+
141
+ In some places, the opportunity is taken when roads are being repaved to lay tramlines (though without erecting overhead cables) even though no service is immediately planned: such is the case in Leipzigerstraße in Berlin, the Haarlemmer Houttuinen in Amsterdam, and Botermarkt in Ghent.
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+
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+ Tram systems operate across national borders in Basel (from Switzerland into France and Germany) and Strasbourg (From France into Germany). It is planned to open a line linking Hasselt (Belgium) with Maastricht (Netherlands) in 2021.
144
+
145
+ Tramway track can have different rail profiles to accommodate the various operating environments of the vehicle. They may be embedded into concrete for street-running operation, or use standard ballasted track with railroad ties on high-speed sections. A more ecological solution is to embed tracks into grass turf.
146
+
147
+ Tramway tracks use a grooved rail with a groove designed for tramway or railway track in pavement or grassed surfaces (grassed track or track in a lawn). The rail has the railhead on one side and the guard on the other. The guard provides accommodation for the flange. The guard carries no weight, but may act as a checkrail. Grooved rail was invented in 1852 by Alphonse Loubat, a French inventor who developed improvements in tram and rail equipment, and helped develop tram lines in New York City and Paris. The invention of grooved rail enabled tramways to be laid without causing a nuisance to other road users, except unsuspecting cyclists, who could get their wheels caught in the groove. The grooves may become filled with gravel and dirt (particularly if infrequently used or after a period of idleness) and need clearing from time to time, this being done by a "scrubber" tram. Failure to clear the grooves can lead to a bumpy ride for the passengers, damage to either wheel or rail and possibly derailing.
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+ In narrow situations double-track tram lines sometimes reduce to single track, or, to avoid switches, have the tracks interlaced, e.g. in the Leidsestraat in Amsterdam on three short stretches (see map detail); this is known as interlaced or gauntlet track. There is a UK example of interlaced track on the Tramlink, just west of Mitcham Station, where the formation is narrowed by an old landslip causing an obstruction. (See photo in Tramlink entry).
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+
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+ Historically, the track gauge has had considerable variations, with narrow gauge common in many early systems. However, most light rail systems are now standard gauge. An important advantage of standard gauge is that standard railway maintenance equipment can be used on it, rather than custom-built machinery. Using standard gauge also allows light rail vehicles to be delivered and relocated conveniently using freight railways and locomotives.
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+ Another factor favoring standard gauge is that low-floor vehicles are becoming popular, and there is generally insufficient space for wheelchairs to move between the wheels in a narrow gauge layout. Standard gauge also enables – at least in theory – a larger choice of manufacturers and thus lower procurement costs for new vehicles. However, other factors such as electrification or loading gauge for which there is more variation may require costly custom built units regardless.
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+ Tram stops may be similar to bus stops in design and use, particularly in street-running sections, where in some cases other vehicles are legally required to stop clear of the tram doors. Some stops may resemble to railway platforms, particularly in private right-of-way sections and where trams are boarded at standard railway platform height, as opposed to using steps at the doorway or low-floor trams.
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+ Approximately 5,000 new trams are manufactured each year. As of February 2017, 4,478 new trams were on order from their makers, with options being open for a further 1,092.[76]
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+ The main manufacturers are:
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+ Trams are in a period of growth, with about 800 tram systems operating around the world, 10 or so new systems being opened each year, and many being gradually extended.[87] Some of these systems date from the late 19th or early 20th centuries. In the past 20 years their numbers have been augmented by modern tramway or light rail systems in cities that had discarded this form of transport. There have also been some new tram systems in cities that never previously had them.
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+ Tramways with tramcars (British English) or street railways with streetcars (North American English) were common throughout the industrialised world in the late 19th and early 20th centuries but they had disappeared from most British, Canadian, French and US cities by the mid-20th century.[88]
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+ By contrast, trams in parts of continental Europe continued to be used by many cities, although there were contractions in some countries, including the Netherlands.[89]
166
+
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+ Since 1980 trams have returned to favour in many places, partly because their tendency to dominate the roadway, formerly seen as a disadvantage, is now considered to be a merit since it raises the visibility of public transport (encouraging car users to change their mode of travel), and enables streets to be reconfigured to give more space to pedestrians, making cites more pleasant places to live. New systems have been built in the United States, United Kingdom, Ireland, Italy, France, Australia and many other countries.
168
+
169
+ In Milan, Italy, the old "Ventotto" trams are considered by its inhabitants a "symbol" of the city. The same can be said of trams in Melbourne in general, but particularly the iconic W class. The Toronto streetcar system had similarly become an iconic symbol of the city, operating the largest network in the Americas as well as the only large-scale tram system in Canada (not including light rail systems, or heritage lines).[90][91]
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+
171
+ The largest tram ((classic tram, streetcar, straßenbahn) and fast tram (light rail, stadtbahn)) networks in the world by route length (as of 2016)[92] are:
172
+
173
+ Other large transit networks that operate streetcar and light rail systems include:
174
+
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+ This list is not exhaustive.
176
+
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+
178
+
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+ Historically, the Paris Tram System was, at its peak, the world's largest system, with 1,111 km (690 mi) of track in 1925[citation needed] (according to other sources, ca. 640 km (400 mi) of route length in 1930). However it was completely closed in 1938.[133] The next largest system appears to have been 857 km (533 mi), in Buenos Aires before 19 February 1963. The third largest was Chicago, with over 850 km (530 mi) of track,[134] but it was all converted to trolleybus and bus services by 21 June 1958. Before its decline, the BVG in Berlin operated a very large network with 634 km (394 mi) of route. Before its system started to be converted to trolleybus (and later bus) services in the 1930s (last tramway closed 6 July 1952), the first-generation London network had 555 km (345 mi) of route in 1931.[135] In 1958 trams in Rio de Jainero were employed on (433 km; 269 mi) of track. The final line, the Santa teresa route was closed in 1968.[136] During a period in the 1980s, the world's largest tram system was in Leningrad (now known as St. Petersburg) with 350 km (220 mi), USSR, and was included as such in the Guinness World Records;[citation needed] however Saint Petersburg's tram system has declined in size since the fall of the Soviet Union. Vienna in 1960 had 340 km (211 mi), before the expansion of bus services and the opening of a subway (1976). Substituting subway services for tram routes continues. 320 km (199 mi) was in Minneapolis-Saint Paul in 1947: There streetcars ended 31 October 1953 in Minneapolis and 19 June 1954 in St. Paul.[137] The Sydney tram network, before it was closed on 25 February 1961, had 291 km (181 mi) of route, and was thus the largest in Australia. As from 1961, the Melbourne system (currently recognised as the world's largest) took over Sydney's title as the largest network in Australia.
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+
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+ In many European cities, much tramway infrastructure was lost in the mid-20th century, though not always on the same scale as in other parts of the world such as North America. Most of Central and Eastern Europe retained the majority of its tramway systems and it is here that the largest and busiest tram systems in the world are found.
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+ Whereas most systems and vehicles in the tram sector are found in Central and Eastern Europe, in the 1960s and 1970s, tram systems were shut down in many places in Western Europe, however urban transportation has been experiencing a sustained long running revival since the 1990s. Many European cities are rehabilitating, upgrading, expanding and reconstructing their old tramway lines and building new tramway lines.[139]
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+
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+ In North America, these vehicles are called "streetcars" (or "trolleys"); the term tram is more likely to be understood as an aerial tramway or a people-mover. Streetcar systems were developed in late 19th to early 20th centuries in a number of cities throughout North America. However, most North American cities saw its streetcar lines removed in the mid-20th century for a variety of financial, technological and social reasons. Exceptions included Boston,[140] Cleveland, Mexico City, New Orleans, Newark, Philadelphia, Pittsburgh, San Francisco, and Toronto.
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+
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+ Toronto currently operates the largest streetcar system in the Americas in terms of track length and ridership. Operated by the Toronto Transit Commission, the streetcar system is the only large-scale streetcar system existing in Canada, excluding heritage streetcar, or light rail systems that are operated in other Canadian municipalities. The streetcar system was established in 1861, and used a variety of vehicles in its history, including horse-drawn streetcars, Peter Witt streetcars, the PCC streetcar, and the Canadian Light Rail Vehicle and its articulated counterpart, the Articulated Light Rail Vehicle. Since December 29, 2019,[141] the system exclusively uses the Flexity Outlook made by Bombardier Transportation.[142][143][144][145]
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+
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+ Streetcars once existed in the Canadian cities of Calgary, Edmonton, Halifax, Hamilton, Kingston, Kitchener, London, Montreal, Ottawa, Peterborough, Quebec City, Regina, Saskatoon, Windsor, and Vancouver. However, Canadian cities excluding Toronto, removed their streetcar systems in the mid-20th century. In the late 1970s and early 1980s, light rail systems were introduced in Calgary and Edmonton; with another light rail system established in Ottawa in 2001. There is now something of a renaissance for light railways in mid-sized cities with Waterloo, Ontario the first to come on line and construction underway in Mississauga, Ontario. In the late 20th century, several Canadian locales restored portions of their defunct streetcar lines, operating them as a heritage feature for tourists. Heritage streetcar lines in Canada include the High Level Bridge Streetcar in Edmonton, the Nelson Electric Tramway in Nelson, and the Whitehorse Waterfront Trolley in Whitehorse.
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+ Pittsburgh had kept most of its streetcar system serving the city and many suburbs, making it the longest-lasting large-network streetcar system in the United States.[citation needed] However, most of Pittsburgh's surviving streetcar lines were converted to light rail in the 1980s. San Francisco's Muni Metro system is the largest surviving streetcar system in the United States, and has even revived previously closed streetcar lines such as the F Market & Wharves heritage streetcar line. In the late 20th century, several cities installed modern light rail systems, in part along the same corridors as their old streetcars systems, the first of these being the San Diego Trolley in San Diego in 1981.
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+
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+ In the 1980s, some cities in the United States brought back streetcars lines, including Memphis, Tampa, and Little Rock; However, these streetcar systems were designed as heritage streetcar lines, and used vintage or replica-vintage vehicles. The first "second-generation streetcar systems" in North America was opened in Portland in 2001.[146] The "second-generation streetcar system," utilizes modern vehicles – vehicles that feature low-floor streetcars. These newer streetcar systems were built in several American cities in the early 21st century including Atlanta, Charlotte, Cincinnati, Dallas, Detroit, Kansas City, Milwaukee, Oklahoma City, Seattle, Tucson, and Washington, D.C..
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+ Model trams are popular in HO scale (1:87) and O scale (1:48 in the US and generally 1:43,5 and 1:45 in Europe and Asia). They are typically powered and will accept plastic figures inside. Common manufacturers are Roco and Lima, with many custom models being made as well. The German firm Hödl[165] and the Austrian Halling[166] specialise in 1:87 scale.[167]
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+
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+ In the US, Bachmann Industries is a mass supplier of HO streetcars and kits. Bowser Manufacturing has produced white metal models for over 50 years.[168] There are many boutique vendors offering limited run epoxy and wood models. At the high end are highly detailed brass models which are usually imported from Japan or Korea and can cost in excess of $500. Many of these run on 16.5 mm (0.65 in) gauge track, which is correct for the representation of 4 ft 8 1⁄2 in (1,435 mm) (standard gauge) in HO scale as in US and Japan, but incorrect in 4 mm (1:76.2) scale, as it represents 4 ft 8 1⁄2 in (1,435 mm). This scale/gauge hybrid is called OO scale.
198
+ O scale trams are also very popular among tram modellers because the increased size allows for more detail and easier crafting of overhead wiring. In the US these models are usually purchased in epoxy or wood kits and some as brass models. The Saint Petersburg Tram Company[169] produces highly detailed polyurethane non-powered O Scale models from around the world which can easily be powered by trucks from vendors like Q-Car.[170]
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+ In the US, one of the best resources for model tram enthusiasts is the East Penn Traction Club of Philadelphia [171] and Trolleyville a website of the Southern California Traction Club.[172]
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+ It is thought that the first example of a working model tramcar in the UK built by an amateur for fun was in 1929, when Frank E. Wilson created a replica of London County Council Tramways E class car 444 in 1:16 scale, which he demonstrated at an early Model Engineer Exhibition. Another of his models was London E/1 1800, which was the only tramway exhibit in the Faraday Memorial Exhibition of 1931. Together with likeminded friends, Frank Wilson went on to found the Tramway & Light Railway Society[173] in 1938, establishing tramway modelling as a hobby.
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+ Mr. Nathan was a passenger by No. 2 tramway car [...] [he] alighted from the car at the southern end, but before he got clear of the rails the car moved onwards [...] he was thus whirled round by the sudden motion of the carriage and his body was brought under the front wheel.
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+ Digestion is the breakdown of large insoluble food molecules into small water-soluble food molecules so that they can be absorbed into the watery blood plasma. In certain organisms, these smaller substances are absorbed through the small intestine into the blood stream. Digestion is a form of catabolism that is often divided into two processes based on how food is broken down: mechanical and chemical digestion. The term mechanical digestion refers to the physical breakdown of large pieces of food into smaller pieces which can subsequently be accessed by digestive enzymes. In chemical digestion, enzymes break down food into the small molecules the body can use.
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+ In the human digestive system, food enters the mouth and mechanical digestion of the food starts by the action of mastication (chewing), a form of mechanical digestion, and the wetting contact of saliva. Saliva, a liquid secreted by the salivary glands, contains salivary amylase, an enzyme which starts the digestion of starch in the food; the saliva also contains mucus, which lubricates the food, and hydrogen carbonate, which provides the ideal conditions of pH (alkaline) for amylase to work. After undergoing mastication and starch digestion, the food will be in the form of a small, round slurry mass called a bolus. It will then travel down the esophagus and into the stomach by the action of peristalsis. Gastric juice in the stomach starts protein digestion. Gastric juice mainly contains hydrochloric acid and pepsin. In infants and toddlers gastric juice also contains rennin. As the first two chemicals may damage the stomach wall, mucus is secreted by the stomach, providing a slimy layer that acts as a shield against the damaging effects of the chemicals. At the same time protein digestion is occurring, mechanical mixing occurs by peristalsis, which is waves of muscular contractions that move along the stomach wall. This allows the mass of food to further mix with the digestive enzymes.
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+ After some time (typically 1–2 hours in humans, 4–6 hours in dogs, 3–4 hours in house cats),[citation needed] the resulting thick liquid is called chyme. When the pyloric sphincter valve opens, chyme enters the duodenum where it mixes with digestive enzymes from the pancreas and bile juice from the liver and then passes through the small intestine, in which digestion continues. When the chyme is fully digested, it is absorbed into the blood. 95% of nutrient absorption occurs in the small intestine. Water and minerals are reabsorbed back into the blood in the colon (large intestine) where the pH is slightly acidic about 5.6 ~ 6.9. Some vitamins, such as biotin and vitamin K (K2MK7) produced by bacteria in the colon are also absorbed into the blood in the colon. Waste material is eliminated from the rectum during defecation.[1]
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+ Digestive systems take many forms. There is a fundamental distinction between internal and external digestion. External digestion developed earlier in evolutionary history, and most fungi still rely on it.[2] In this process, enzymes are secreted into the environment surrounding the organism, where they break down an organic material, and some of the products diffuse back to the organism. Animals have a tube (gastrointestinal tract) in which internal digestion occurs, which is more efficient because more of the broken down products can be captured, and the internal chemical environment can be more efficiently controlled.[3]
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+ Some organisms, including nearly all spiders, simply secrete biotoxins and digestive chemicals (e.g., enzymes) into the extracellular environment prior to ingestion of the consequent "soup". In others, once potential nutrients or food is inside the organism, digestion can be conducted to a vesicle or a sac-like structure, through a tube, or through several specialized organs aimed at making the absorption of nutrients more efficient.
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+ Bacteria use several systems to obtain nutrients from other organisms in the environments.
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+ In a channel transupport system, several proteins form a contiguous channel traversing the inner and outer membranes of the bacteria. It is a simple system, which consists of only three protein subunits: the ABC protein, membrane fusion protein (MFP), and outer membrane protein (OMP)[specify]. This secretion system transports various molecules, from ions, drugs, to proteins of various sizes (20–900 kDa). The molecules secreted vary in size from the small Escherichia coli peptide colicin V, (10 kDa) to the Pseudomonas fluorescens cell adhesion protein LapA of 900 kDa.[4]
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+ A type III secretion system means that a molecular syringe is used through which a bacterium (e.g. certain types of Salmonella, Shigella, Yersinia) can inject nutrients into protist cells. One such mechanism was first discovered in Y. pestis and showed that toxins could be injected directly from the bacterial cytoplasm into the cytoplasm of its host's cells rather than simply be secreted into the extracellular medium.[5]
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+ The conjugation machinery of some bacteria (and archaeal flagella) is capable of transporting both DNA and proteins. It was discovered in Agrobacterium tumefaciens, which uses this system to introduce the Ti plasmid and proteins into the host, which develops the crown gall (tumor).[6] The VirB complex of Agrobacterium tumefaciens is the prototypic system.[7]
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+ The nitrogen fixing Rhizobia are an interesting case, wherein conjugative elements naturally engage in inter-kingdom conjugation. Such elements as the Agrobacterium Ti or Ri plasmids contain elements that can transfer to plant cells. Transferred genes enter the plant cell nucleus and effectively transform the plant cells into factories for the production of opines, which the bacteria use as carbon and energy sources. Infected plant cells form crown gall or root tumors. The Ti and Ri plasmids are thus endosymbionts of the bacteria, which are in turn endosymbionts (or parasites) of the infected plant.
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+ The Ti and Ri plasmids are themselves conjugative. Ti and Ri transfer between bacteria uses an independent system (the tra, or transfer, operon) from that for inter-kingdom transfer (the vir, or virulence, operon). Such transfer creates virulent strains from previously avirulent Agrobacteria.
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+ In addition to the use of the multiprotein complexes listed above, Gram-negative bacteria possess another method for release of material: the formation of outer membrane vesicles.[8][9] Portions of the outer membrane pinch off, forming spherical structures made of a lipid bilayer enclosing periplasmic materials. Vesicles from a number of bacterial species have been found to contain virulence factors, some have immunomodulatory effects, and some can directly adhere to and intoxicate host cells. While release of vesicles has been demonstrated as a general response to stress conditions, the process of loading cargo proteins seems to be selective.[10]
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+ The gastrovascular cavity functions as a stomach in both digestion and the distribution of nutrients to all parts of the body. Extracellular digestion takes place within this central cavity, which is lined with the gastrodermis, the internal layer of epithelium. This cavity has only one opening to the outside that functions as both a mouth and an anus: waste and undigested matter is excreted through the mouth/anus, which can be described as an incomplete gut.
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+ In a plant such as the Venus Flytrap that can make its own food through photosynthesis, it does not eat and digest its prey for the traditional objectives of harvesting energy and carbon, but mines prey primarily for essential nutrients (nitrogen and phosphorus in particular) that are in short supply in its boggy, acidic habitat.[11]
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+ A phagosome is a vacuole formed around a particle absorbed by phagocytosis. The vacuole is formed by the fusion of the cell membrane around the particle. A phagosome is a cellular compartment in which pathogenic microorganisms can be killed and digested. Phagosomes fuse with lysosomes in their maturation process, forming phagolysosomes. In humans, Entamoeba histolytica can phagocytose red blood cells.[12]
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+ To aid in the digestion of their food animals evolved organs such as beaks, tongues, teeth, a crop, gizzard, and others.
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+ Birds have bony beaks that are specialised according to the bird's ecological niche. For example, macaws primarily eat seeds, nuts, and fruit, using their impressive beaks to open even the toughest seed. First they scratch a thin line with the sharp point of the beak, then they shear the seed open with the sides of the beak.
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+ The mouth of the squid is equipped with a sharp horny beak mainly made of cross-linked proteins. It is used to kill and tear prey into manageable pieces. The beak is very robust, but does not contain any minerals, unlike the teeth and jaws of many other organisms, including marine species.[13] The beak is the only indigestible part of the squid.
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+ The tongue is skeletal muscle on the floor of the mouth of most vertebrates, that manipulates food for chewing (mastication) and swallowing (deglutition). It is sensitive and kept moist by saliva. The underside of the tongue is covered with a smooth mucous membrane. The tongue also has a touch sense for locating and positioning food particles that require further chewing. The tongue is utilized to roll food particles into a bolus before being transported down the esophagus through peristalsis.
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+ The sublingual region underneath the front of the tongue is a location where the oral mucosa is very thin, and underlain by a plexus of veins. This is an ideal location for introducing certain medications to the body. The sublingual route takes advantage of the highly vascular quality of the oral cavity, and allows for the speedy application of medication into the cardiovascular system, bypassing the gastrointestinal tract.
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+ Teeth (singular tooth) are small whitish structures found in the jaws (or mouths) of many vertebrates that are used to tear, scrape, milk and chew food. Teeth are not made of bone, but rather of tissues of varying density and hardness, such as enamel, dentine and cementum. Human teeth have a blood and nerve supply which enables proprioception. This is the ability of sensation when chewing, for example if we were to bite into something too hard for our teeth, such as a chipped plate mixed in food, our teeth send a message to our brain and we realise that it cannot be chewed, so we stop trying.
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+ The shapes, sizes and numbers of types of animals' teeth are related to their diets. For example, herbivores have a number of molars which are used to grind plant matter, which is difficult to digest. Carnivores have canine teeth which are used to kill and tear meat.
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+ A crop, or croup, is a thin-walled expanded portion of the alimentary tract used for the storage of food prior to digestion. In some birds it is an expanded, muscular pouch near the gullet or throat. In adult doves and pigeons, the crop can produce crop milk to feed newly hatched birds.[14]
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+ Certain insects may have a crop or enlarged esophagus.
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+ Herbivores have evolved cecums (or an abomasum in the case of ruminants). Ruminants have a fore-stomach with four chambers. These are the rumen, reticulum, omasum, and abomasum. In the first two chambers, the rumen and the reticulum, the food is mixed with saliva and separates into layers of solid and liquid material. Solids clump together to form the cud (or bolus). The cud is then regurgitated, chewed slowly to completely mix it with saliva and to break down the particle size.
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+ Fibre, especially cellulose and hemi-cellulose, is primarily broken down into the volatile fatty acids, acetic acid, propionic acid and butyric acid in these chambers (the reticulo-rumen) by microbes: (bacteria, protozoa, and fungi). In the omasum, water and many of the inorganic mineral elements are absorbed into the blood stream.
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+ The abomasum is the fourth and final stomach compartment in ruminants. It is a close equivalent of a monogastric stomach (e.g., those in humans or pigs), and digesta is processed here in much the same way. It serves primarily as a site for acid hydrolysis of microbial and dietary protein, preparing these protein sources for further digestion and absorption in the small intestine. Digesta is finally moved into the small intestine, where the digestion and absorption of nutrients occurs. Microbes produced in the reticulo-rumen are also digested in the small intestine.
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+ Regurgitation has been mentioned above under abomasum and crop, referring to crop milk, a secretion from the lining of the crop of pigeons and doves with which the parents feed their young by regurgitation.[15]
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+ Many sharks have the ability to turn their stomachs inside out and evert it out of their mouths in order to get rid of unwanted contents (perhaps developed as a way to reduce exposure to toxins).
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+ Other animals, such as rabbits and rodents, practise coprophagia behaviours – eating specialised faeces in order to re-digest food, especially in the case of roughage. Capybara, rabbits, hamsters and other related species do not have a complex digestive system as do, for example, ruminants. Instead they extract more nutrition from grass by giving their food a second pass through the gut. Soft faecal pellets of partially digested food are excreted and generally consumed immediately. They also produce normal droppings, which are not eaten.
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+ Young elephants, pandas, koalas, and hippos eat the faeces of their mother, probably to obtain the bacteria required to properly digest vegetation. When they are born, their intestines do not contain these bacteria (they are completely sterile). Without them, they would be unable to get any nutritional value from many plant components.
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+ An earthworm's digestive system consists of a mouth, pharynx, esophagus, crop, gizzard, and intestine. The mouth is surrounded by strong lips, which act like a hand to grab pieces of dead grass, leaves, and weeds, with bits of soil to help chew. The lips break the food down into smaller pieces. In the pharynx, the food is lubricated by mucus secretions for easier passage. The esophagus adds calcium carbonate to neutralize the acids formed by food matter decay. Temporary storage occurs in the crop where food and calcium carbonate are mixed. The powerful muscles of the gizzard churn and mix the mass of food and dirt. When the churning is complete, the glands in the walls of the gizzard add enzymes to the thick paste, which helps chemically breakdown the organic matter. By peristalsis, the mixture is sent to the intestine where friendly bacteria continue chemical breakdown. This releases carbohydrates, protein, fat, and various vitamins and minerals for absorption into the body.
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+ In most vertebrates, digestion is a multistage process in the digestive system, starting from ingestion of raw materials, most often other organisms. Ingestion usually involves some type of mechanical and chemical processing. Digestion is separated into four steps:
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+ Underlying the process is muscle movement throughout the system through swallowing and peristalsis. Each step in digestion requires energy, and thus imposes an "overhead charge" on the energy made available from absorbed substances. Differences in that overhead cost are important influences on lifestyle, behavior, and even physical structures. Examples may be seen in humans, who differ considerably from other hominids (lack of hair, smaller jaws and musculature, different dentition, length of intestines, cooking, etc.).
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+ The major part of digestion takes place in the small intestine. The large intestine primarily serves as a site for fermentation of indigestible matter by gut bacteria and for resorption of water from digests before excretion.
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+ In mammals, preparation for digestion begins with the cephalic phase in which saliva is produced in the mouth and digestive enzymes are produced in the stomach. Mechanical and chemical digestion begin in the mouth where food is chewed, and mixed with saliva to begin enzymatic processing of starches. The stomach continues to break food down mechanically and chemically through churning and mixing with both acids and enzymes. Absorption occurs in the stomach and gastrointestinal tract, and the process finishes with defecation.[1]
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+ The human gastrointestinal tract is around 9 meters long. Food digestion physiology varies between individuals and upon other factors such as the characteristics of the food and size of the meal, and the process of digestion normally takes between 24 and 72 hours.[16]
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+ Digestion begins in the mouth with the secretion of saliva and its digestive enzymes. Food is formed into a bolus by the mechanical mastication and swallowed into the esophagus from where it enters the stomach through the action of peristalsis. Gastric juice contains hydrochloric acid and pepsin which would damage the walls of the stomach and mucus is secreted for protection. In the stomach further release of enzymes break down the food further and this is combined with the churning action of the stomach. The partially digested food enters the duodenum as a thick semi-liquid chyme. In the small intestine, the larger part of digestion takes place and this is helped by the secretions of bile, pancreatic juice and intestinal juice. The intestinal walls are lined with villi, and their epithelial cells is covered with numerous microvilli to improve the absorption of nutrients by increasing the surface area of the intestine.
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+ In the large intestine the passage of food is slower to enable fermentation by the gut flora to take place. Here water is absorbed and waste material stored as feces to be removed by defecation via the anal canal and anus.
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+ Different phases of digestion take place including: the cephalic phase, gastric phase, and intestinal phase.
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+ The cephalic phase occurs at the sight, thought and smell of food, which stimulate the cerebral cortex. Taste and smell stimuli are sent to the hypothalamus and medulla oblongata. After this it is routed through the vagus nerve and release of acetylcholine. Gastric secretion at this phase rises to 40% of maximum rate. Acidity in the stomach is not buffered by food at this point and thus acts to inhibit parietal (secretes acid) and G cell (secretes gastrin) activity via D cell secretion of somatostatin.
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+
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+ The gastric phase takes 3 to 4 hours. It is stimulated by distension of the stomach, presence of food in stomach and decrease in pH. Distention activates long and myenteric reflexes. This activates the release of acetylcholine, which stimulates the release of more gastric juices. As protein enters the stomach, it binds to hydrogen ions, which raises the pH of the stomach. Inhibition of gastrin and gastric acid secretion is lifted. This triggers G cells to release gastrin, which in turn stimulates parietal cells to secrete gastric acid. Gastric acid is about 0.5% hydrochloric acid (HCl), which lowers the pH to the desired pH of 1–3. Acid release is also triggered by acetylcholine and histamine.
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+ The intestinal phase has two parts, the excitatory and the inhibitory. Partially digested food fills the duodenum. This triggers intestinal gastrin to be released. Enterogastric reflex inhibits vagal nuclei, activating sympathetic fibers causing the pyloric sphincter to tighten to prevent more food from entering, and inhibits local reflexes.
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+ Protein digestion occurs in the stomach and duodenum in which 3 main enzymes, pepsin secreted by the stomach and trypsin and chymotrypsin secreted by the pancreas, break down food proteins into polypeptides that are then broken down by various exopeptidases and dipeptidases into amino acids. The digestive enzymes however are mostly secreted as their inactive precursors, the zymogens. For example, trypsin is secreted by pancreas in the form of trypsinogen, which is activated in the duodenum by enterokinase to form trypsin. Trypsin then cleaves proteins to smaller polypeptides.
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+
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+ Digestion of some fats can begin in the mouth where lingual lipase breaks down some short chain lipids into diglycerides. However fats are mainly digested in the small intestine.[17] The presence of fat in the small intestine produces hormones that stimulate the release of pancreatic lipase from the pancreas and bile from the liver which helps in the emulsification of fats for absorption of fatty acids.[17] Complete digestion of one molecule of fat (a triglyceride) results a mixture of fatty acids, mono- and di-glycerides, as well as some undigested triglycerides, but no free glycerol molecules.[17]
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+
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+ In humans, dietary starches are composed of glucose units arranged in long chains called amylose, a polysaccharide. During digestion, bonds between glucose molecules are broken by salivary and pancreatic amylase, resulting in progressively smaller chains of glucose. This results in simple sugars glucose and maltose (2 glucose molecules) that can be absorbed by the small intestine.
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+
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+ Lactase is an enzyme that breaks down the disaccharide lactose to its component parts, glucose and galactose. Glucose and galactose can be absorbed by the small intestine. Approximately 65 percent of the adult population produce only small amounts of lactase and are unable to eat unfermented milk-based foods. This is commonly known as lactose intolerance. Lactose intolerance varies widely by genetic heritage; more than 90 percent of peoples of east Asian descent are lactose intolerant, in contrast to about 5 percent of people of northern European descent.[18]
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+ Sucrase is an enzyme that breaks down the disaccharide sucrose, commonly known as table sugar, cane sugar, or beet sugar. Sucrose digestion yields the sugars fructose and glucose which are readily absorbed by the small intestine.
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+
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+ DNA and RNA are broken down into mononucleotides by the nucleases deoxyribonuclease and ribonuclease (DNase and RNase) from the pancreas.
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+
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+ Some nutrients are complex molecules (for example vitamin B12) which would be destroyed if they were broken down into their functional groups. To digest vitamin B12 non-destructively, haptocorrin in saliva strongly binds and protects the B12 molecules from stomach acid as they enter the stomach and are cleaved from their protein complexes.[19]
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+
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+ After the B12-haptocorrin complexes pass from the stomach via the pylorus to the duodenum, pancreatic proteases cleave haptocorrin from the B12 molecules which rebind to intrinsic factor (IF). These B12-IF complexes travel to the ileum portion of the small intestine where cubilin receptors enable assimilation and circulation of B12-IF complexes in the blood.[20]
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+
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+ There are at least five hormones that aid and regulate the digestive system in mammals. There are variations across the vertebrates, as for instance in birds. Arrangements are complex and additional details are regularly discovered. For instance, more connections to metabolic control (largely the glucose-insulin system) have been uncovered in recent years.
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+
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+ Digestion is a complex process controlled by several factors. pH plays a crucial role in a normally functioning digestive tract. In the mouth, pharynx and esophagus, pH is typically about 6.8, very weakly acidic. Saliva controls pH in this region of the digestive tract. Salivary amylase is contained in saliva and starts the breakdown of carbohydrates into monosaccharides. Most digestive enzymes are sensitive to pH and will denature in a high or low pH environment.
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+ The stomach's high acidity inhibits the breakdown of carbohydrates within it. This acidity confers two benefits: it denatures proteins for further digestion in the small intestines, and provides non-specific immunity, damaging or eliminating various pathogens.[citation needed]
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+
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+ In the small intestines, the duodenum provides critical pH balancing to activate digestive enzymes. The liver secretes bile into the duodenum to neutralize the acidic conditions from the stomach, and the pancreatic duct empties into the duodenum, adding bicarbonate to neutralize the acidic chyme, thus creating a neutral environment. The mucosal tissue of the small intestines is alkaline with a pH of about 8.5.[citation needed]
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1
+
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+
3
+ A tram (in North America streetcar or trolley) is a rail vehicle that runs on tramway tracks along public urban streets; some include segments of segregated right-of-way.[1][2] The lines or networks operated by tramcars are called tramways. Historically the term electric street railways was also used in the United States. In the United States, the term tram has sometimes been used for rubber-tired trackless trains, which are unrelated to other kinds of trams.
4
+
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+ Tram vehicles are usually lighter and shorter than main line and rapid transit trains. Today, most trams use electrical power, usually fed by a pantograph sliding on an overhead line; older systems may use a trolley pole or a bow collector. In some cases, a contact shoe on a third rail is used. If necessary, they may have dual power systems—electricity in city streets and diesel in more rural environments. Occasionally, trams also carry freight.
6
+
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+ Trams are now commonly included in the wider term "light rail",[3] which also includes grade-separated systems. Some trams, known as tram-trains, may have segments that run on mainline railway tracks, similar to interurban systems. The differences between these modes of rail transport are often indistinct and a given system may combine multiple features.
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+
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+ One of the advantages over earlier forms of transit was the low rolling resistance of metal wheels on steel rails, allowing the trams to haul a greater load for a given effort. Problems included the high total cost of ownership of horses. Electric trams largely replaced animal power in the late 19th and early 20th centuries. Improvements in other vehicles such as buses led to decline of trams in the mid 20th century. However, trams have seen resurgence in recent years.
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+
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+ The English terms tram and tramway are derived from the Scots word tram,[4] referring respectively to a type of truck (goods wagon or freight railroad car) used in coal mines and the tracks on which they ran. The word tram probably derived from Middle Flemish trame ("beam, handle of a barrow, bar, rung"). The identical word la trame with the meaning "crossbeam" is also used in the French language. Etymologists believe that the word tram refers to the wooden beams the railway tracks were initially made of before the railroad pioneers switched to the much more wear-resistant tracks made of iron and, later, steel.[5] The word Tram-car is attested from 1873.[6]
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+
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+ Although the terms tram and tramway have been adopted by many languages, they are not used universally in English; North Americans prefer streetcar, trolley, or trolleycar. The term streetcar is first recorded in 1840, and originally referred to horsecars. When electrification came, Americans began to speak of trolleycars or later, trolleys. A widely held belief holds the word to derive from the troller (said to derive from the words traveler and roller), a four-wheeled device that was dragged along dual overhead wires by a cable that connected the troller to the top of the car and collected electrical power from the overhead wires;[7] this portmanteau derivation is, however, most likely folk etymology. "Trolley" and variants refer to the verb troll, meaning "roll" and probably derived from Old French,[8] and cognate uses of the word were well established for handcarts and horse drayage, as well as for nautical uses.[9]
14
+
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+ The alternative North American term 'trolley' may strictly speaking be considered incorrect, as the term can also be applied to cable cars, or conduit cars that instead draw power from an underground supply. Conventional diesel tourist buses decorated to look like streetcars are sometimes called trolleys in the US (tourist trolley). Furthering confusion, the term tram has instead been applied to open-sided, low-speed segmented vehicles on rubber tires generally used to ferry tourists short distances, for example on the Universal Studios backlot tour and, in many countries, as tourist transport to major destinations. The term may also apply to an aerial ropeway, e.g. the Roosevelt Island Tramway.
16
+
17
+ Although the use of the term trolley for tram was not adopted in Europe, the term was later associated with the trolleybus, a rubber-tyred vehicle running on hard pavement, which draws its power from pairs of overhead wires. These electric buses, which use twin trolley poles, are also called trackless trolleys (particularly in the northeastern US), or sometimes simply trolleys (in the UK, as well as the Pacific Northwest, including Seattle, and Vancouver).
18
+
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+ The New South Wales government in Australia has decided to use the term "light rail" for their trams.
20
+
21
+ The history of trams, streetcars or trolley systems, began in the early nineteenth century. It can be divided up into several discrete periods defined by the principal means of motive power used.
22
+
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+ The world's first passenger train or tram was the Swansea and Mumbles Railway, in Wales, UK. The Mumbles Railway Act was passed by the British Parliament in 1804, and horse-drawn service started in 1807.[10] The service closed in 1827, but was restarted in 1860, again using horses.[11] It was worked by steam from 1877, and then, from 1929, by very large (106-seater) electric tramcars, until closure in 1961.[citation needed] The Swansea and Mumbles Railway was something of a one-off however, and no street tramway would appear in Britain until 1860 when one was built in Birkenhead by the American George Francis Train.[12]
24
+
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+ Street railways developed in America before Europe, largely due to the poor paving of the streets in American cities which made them unsuitable for horsebuses, which were then common on the well-paved streets of European cities. Running the horsecars on rails allowed for a much smoother ride. There are records of a street railway running in Baltimore as early as 1828, however the first authenticated streetcar in America, was the New York and Harlem Railroad developed by the Irish coach builder John Stephenson, in New York City which began service in the year 1832.[13][14] The New York and Harlem Railroad's Fourth Avenue Line ran along the Bowery and Fourth Avenue in New York City. It was followed in 1835 by the New Orleans and Carrollton Railroad in New Orleans, Louisiana,[15] which still operates as the St. Charles Streetcar Line. Other American cities did not follow until the 1850s, after which the "animal railway" became an increasingly common feature in the larger towns.[15]
26
+
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+ The first permanent tram line in continental Europe was opened in Paris in 1855 by Alphonse Loubat who had previously worked on American streetcar lines.[16] The tram was developed in numerous cities of Europe (some of the most extensive systems were found in Berlin, Budapest, Birmingham, Leningrad, Lisbon, London, Manchester, Paris).
28
+
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+ The first tram in South America opened in 1858 in Santiago, Chile. The first trams in Australia opened in 1860 in Sydney. Africa's first tram service started in Alexandria on 8 January 1863. The first trams in Asia opened in 1869 in Batavia (now Jakarta), Netherlands East Indies (now Indonesia).
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+
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+ Problems with horsecars included the fact that any given animal could only work so many hours on a given day, had to be housed, groomed, fed and cared for day in and day out, and produced prodigious amounts of manure, which the streetcar company was charged with storing and then disposing of. Since a typical horse pulled a streetcar for about a dozen miles a day and worked for four or five hours, many systems needed ten or more horses in stable for each horsecar.
32
+
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+ Horsecars were largely replaced by electric-powered trams following the improvement of an overhead trolley system on trams for collecting electricity from overhead wires by Frank J. Sprague. His spring-loaded trolley pole used a wheel to travel along the wire. In late 1887 and early 1888, using his trolley system, Sprague installed the first successful large electric street railway system in Richmond, Virginia. Within a year, the economy of electric power had replaced more costly horsecars in many cities. By 1889, 110 electric railways incorporating Sprague's equipment had been begun or planned on several continents.[17]
34
+
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+ Horses continued to be used for light shunting well into the 20th century, and many large metropolitan lines lasted into the early 20th century. New York City had a regular horsecar service on the Bleecker Street Line until its closure in 1917.[18] Pittsburgh, Pennsylvania, had its Sarah Street line drawn by horses until 1923. The last regular mule-drawn cars in the US ran in Sulphur Rock, Arkansas, until 1926 and were commemorated by a U.S. postage stamp issued in 1983.[19] The last mule tram service in Mexico City ended in 1932, and a mule tram in Celaya, Mexico, survived until 1954.[20] The last horse-drawn tram to be withdrawn from public service in the UK took passengers from Fintona railway station to Fintona Junction one mile away on the main Omagh to Enniskillen railway in Northern Ireland. The tram made its last journey on 30 September 1957 when the Omagh to Enniskillen line closed. The "van" now lies at the Ulster Transport Museum.
36
+
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+ Horse-drawn trams still operate on the 1876-built Douglas Bay Horse Tramway in the Isle of Man, and at the 1894-built horse tram at Victor Harbor in South Australia. New horse-drawn systems have been established at the Hokkaidō Museum in Japan and also in Disneyland. A horse tram route in Polish gmina Mrozy, first built in 1902, was reopened in 2012.
38
+
39
+ The first mechanical trams were powered by steam. Generally, there were two types of steam tram. The first and most common had a small steam locomotive (called a tram engine in the UK) at the head of a line of one or more carriages, similar to a small train. Systems with such steam trams included Christchurch, New Zealand; Sydney, Australia; other city systems in New South Wales; Munich, Germany (from August 1883 on),[21] British India (Pakistan) (from 1885) and the Dublin & Blessington Steam Tramway (from 1888) in Ireland. Steam tramways also were used on the suburban tramway lines around Milan and Padua; the last Gamba de Legn ("Peg-Leg") tramway ran on the Milan-Magenta-Castano Primo route in late 1957.[22]
40
+
41
+ The other style of steam tram had the steam engine in the body of the tram, referred to as a tram engine (UK) or steam dummy (US). The most notable system to adopt such trams was in Paris. French-designed steam trams also operated in Rockhampton, in the Australian state of Queensland between 1909 and 1939. Stockholm, Sweden, had a steam tram line at the island of Södermalm between 1887 and 1901.
42
+
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+ Tram engines usually had modifications to make them suitable for street running in residential areas. The wheels, and other moving parts of the machinery, were usually enclosed for safety reasons and to make the engines quieter. Measures were often taken to prevent the engines from emitting visible smoke or steam. Usually the engines used coke rather than coal as fuel to avoid emitting smoke; condensers or superheating were used to avoid emitting visible steam. A major drawback of this style of tram was the limited space for the engine, so that these trams were usually underpowered. Steam tram engines faded out around the 1890s to 1900s, being replaced by electric trams.
44
+
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+ Another motive system for trams was the cable car, which was pulled along a fixed track by a moving steel cable. The power to move the cable was normally provided at a "powerhouse" site a distance away from the actual vehicle. The London and Blackwall Railway, which opened for passengers in east London, England, in 1840 used such a system.[23]
46
+
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+ The first practical cable car line was tested in San Francisco, in 1873. Part of its success is attributed to the development of an effective and reliable cable grip mechanism, to grab and release the moving cable without damage. The second city to operate cable trams was Dunedin in New Zealand, from 1881 to 1957.
48
+
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+ The most extensive cable system in the US was built in Chicago, having been built in stages between 1859 and 1892. New York City developed multiple cable car lines, that operated from 1883 to 1909.[24] Los Angeles also had several cable car lines, including the Second Street Cable Railroad, which operated from 1885 to 1889, and the Temple Street Cable Railway, which operated from 1886 to 1898.
50
+
51
+ From 1885 to 1940, the city of Melbourne, Victoria, Australia operated one of the largest cable systems in the world, at its peak running 592 trams on 75 kilometres (47 mi) of track. There were also two isolated cable lines in Sydney, New South Wales, Australia; the North Sydney line from 1886 to 1900,[25] and the King Street line from 1892 to 1905.
52
+
53
+ In Dresden, Germany, in 1901 an elevated suspended cable car following the Eugen Langen one-railed floating tram system started operating. Cable cars operated on Highgate Hill in North London and Kennington to Brixton Hill In South London.[when?] They also worked around "Upper Douglas" in the Isle of Man from 1897 to 1929 (cable car 72/73 is the sole survivor of the fleet).
54
+
55
+ Cable cars suffered from high infrastructure costs, since an expensive system of cables, pulleys, stationary engines and lengthy underground vault structures beneath the rails had to be provided. They also required physical strength and skill to operate, and alert operators to avoid obstructions and other cable cars. The cable had to be disconnected ("dropped") at designated locations to allow the cars to coast by inertia, for example when crossing another cable line. The cable would then have to be "picked up" to resume progress, the whole operation requiring precise timing to avoid damage to the cable and the grip mechanism. Breaks and frays in the cable, which occurred frequently, required the complete cessation of services over a cable route while the cable was repaired. Due to overall wear, the entire length of cable (typically several kilometres) would have to be replaced on a regular schedule. After the development of reliable electrically powered trams, the costly high-maintenance cable car systems were rapidly replaced in most locations.
56
+
57
+ Cable cars remained especially effective in hilly cities, since their nondriven wheels would not lose traction as they climbed or descended a steep hill. The moving cable would physically pull the car up the hill at a steady pace, unlike a low-powered steam or horse-drawn car. Cable cars do have wheel brakes and track brakes, but the cable also helps restrain the car to going downhill at a constant speed. Performance in steep terrain partially explains the survival of cable cars in San Francisco.
58
+
59
+ The San Francisco cable cars, though significantly reduced in number, continue to perform a regular transportation function, in addition to being a well-known tourist attraction. A single cable line also survives in Wellington, New Zealand (rebuilt in 1979 as a funicular but still called the "Wellington Cable Car"). Another system, actually two separate cable lines with a shared power station in the middle, operates from the Welsh town of Llandudno up to the top of the Great Orme hill in North Wales, UK.
60
+
61
+ In the late 19th and early 20th centuries a number of systems in various parts of the world employed trams powered by gas, naphtha gas or coal gas in particular. Gas trams are known to have operated between Alphington and Clifton Hill in the northern suburbs of Melbourne, Australia (1886–1888); in Berlin and Dresden, Germany; in Estonia (1921–1951); between Jelenia Góra, Cieplice, and Sobieszów in Poland (from 1897); and in the UK at Lytham St Annes, Neath (1896–1920), and Trafford Park, Manchester (1897–1908).
62
+
63
+ On 29 December 1886 the Melbourne newspaper The Argus reprinted a report from the San Francisco Bulletin that Mr Noble had demonstrated a new 'motor car' for tramways 'with success'. The tramcar 'exactly similar in size, shape, and capacity to a cable grip car' had the 'motive power' of gas 'with which the reservoir is to be charged once a day at power stations by means of a rubber hose'. The car also carried an electricity generator for 'lighting up the tram and also for driving the engine on steep grades and effecting a start'.[26]
64
+
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+ Comparatively little has been published about gas trams. However, research on the subject was carried out for an article in the October 2011 edition of "The Times", the historical journal of the Australian Association of Timetable Collectors, now the Australian Timetable Association.[27][28][29][30]
66
+
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+ A tram system powered by compressed natural gas was due to open in Malaysia in 2012,[31] but the news about the project appears to have dried up.
68
+
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+ The world's first electric tram line operated in Sestroretsk near Saint Petersburg, Russia, invented and tested by Fyodor Pirotsky in 1875.[32][33] Later, using a similar technology, Pirotsky put into service the first public electric tramway in St. Petersburg, which operated only during September 1880.[34]
70
+ The second demonstrative tramway was presented by Siemens & Halske at the 1879 Berlin Industrial Exposition.
71
+ The first public electric tramway used for permanent service was the Gross-Lichterfelde tramway in Lichterfelde near Berlin in Germany, which opened in 1881. It was built by Werner von Siemens who contacted Pirotsky. This was world's first commercially successful electric tram. It initially drew current from the rails, with overhead wire being installed in 1883.[35]
72
+
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+ In Britain, Volk's Electric Railway was opened in 1883 in Brighton). This two kilometer line along the seafront, re-gauged to 2 ft 9 in (838 mm) in 1884, remains in service to this day and is the oldest operating electric tramway in the world. Also in 1883, Mödling and Hinterbrühl Tram was opened near Vienna in Austria. It was the first tram in the world in regular service that was run with electricity served by an overhead line with pantograph current collectors. The Blackpool Tramway was opened in Blackpool, UK on 29 September 1885 using conduit collection along Blackpool Promenade. This system is still in operation in a modernised form.[36]
74
+
75
+ Earliest tram system in Canada was by John Joseph Wright, brother of the famous mining entrepreneur Whitaker Wright, in Toronto in 1883, introducing electric trams in 1892. In the US, multiple functioning experimental electric trams were exhibited at the 1884 World Cotton Centennial World's Fair in New Orleans, Louisiana, but they were not deemed good enough to replace the Lamm fireless engines then propelling the St. Charles Avenue Streetcar in that city. The first commercial installation of an electric streetcar in the United States was built in 1884 in Cleveland, Ohio and operated for a period of one year by the East Cleveland Street Railway Company.[37] Trams were operated in Richmond, Virginia, in 1888, on the Richmond Union Passenger Railway built by Frank J. Sprague. Sprague later developed multiple unit control, first demonstrated in Chicago in 1897, allowing multiple cars to be coupled together and operated by a single motorman. This gave rise to the modern subway train. Following the improvement of an overhead "trolley" system on streetcars for collecting electricity from overhead wires by Sprague, electric tram systems were rapidly adopted across the world.[citation needed]
76
+
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+ Earlier electric trains proved difficult or unreliable and experienced limited success until the second half of the 1880s, when new types of current collectors were developed [34]. Siemens' line, for example, provided power through a live rail and a return rail, like a model train, limiting the voltage that could be used, and delivering electric shocks to people and animals crossing the tracks.[38] Siemens later designed his own version of overhead current collection, called the bow collector, and Thorold, Ontario, opened in 1887, and was considered quite successful at the time. While this line proved quite versatile as one of the earliest fully functional electric streetcar installations, it required horse-drawn support while climbing the Niagara Escarpment and for two months of the winter when hydroelectricity was not available. It continued in service in its original form into the 1950s.[citation needed]
78
+
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+ Sidney Howe Short designed and produced the first electric motor that operated a streetcar without gears. The motor had its armature direct-connected to the streetcar's axle for the driving force.[39][40][41][42][43] Short pioneered "use of a conduit system of concealed feed" thereby eliminating the necessity of overhead wire and a trolley pole for street cars and railways.[44][39][40] While at the University of Denver he conducted important experiments which established that multiple unit powered cars were a better way to operate trains and trolleys.[39][40]
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+
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+ Electric tramways spread to many European cities in the 1890s, such as Prague, Bohemia (then in the Austro-Hungarian Empire), in 1891; Kiev, Ukraine, in 1892 (the first permanent electric tram line in the Russian Empire); Dresden, Germany, Lyon, France, and Milan and Genoa, Italy, in 1893; Rome, Italy, Plauen, Germany, in 1894; Bristol, United Kingdom, Munich, in 1895; Bilbao, Spain, in 1896; Copenhagen, Denmark, and Vienna, Austria, in 1897; Florence and Turin, Italy, in 1898; Helsinki, Finland, and Madrid and Barcelona, Spain, in 1899.[34] Sarajevo built a citywide system of electric trams in 1895.[45] Budapest established its tramway system in 1887, and its ring line has grown to be the busiest tram line in Europe, with a tram running every 60 seconds at rush hour. Bucharest and Belgrade[46] ran a regular service from 1894.[47][48] Ljubljana introduced its tram system in 1901 – it closed in 1958.[49] Oslo had the first tramway in Scandinavia, starting operation on 2 March 1894.[50]
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+
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+ The first electric tramway in Australia was a Sprague system demonstrated at the 1888 Melbourne Centennial Exhibition in Melbourne; afterwards, this was installed as a commercial venture operating between the outer Melbourne suburb of Box Hill and the then tourist-oriented country town Doncaster from 1889 to 1896.[51] As well, electric systems were built in Adelaide, Ballarat, Bendigo, Brisbane, Fremantle, Geelong, Hobart, Kalgoorlie, Launceston, Leonora, Newcastle, Perth, and Sydney.
84
+
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+ By the 1970s, the only full tramway system remaining in Australia was the Melbourne tram system. However, there were also a few single lines remaining elsewhere: the Glenelg tram line, connecting Adelaide to the beachside suburb of Glenelg, and tourist trams in the Victorian Goldfields cities of Bendigo and Ballarat. In recent years the Melbourne system, generally recognised as the largest urban tram network in the world, has been considerably modernised and expanded.[52] The Adelaide line has also been extended to the Entertainment Centre, and work is progressing on further extensions.[53] Sydney re-introduced trams (or light rail) on 31 August 1997. A completely new system, known as G:link, was introduced on the Gold Coast, Queensland on 20 July 2014. The Newcastle Light Rail opened in February 2019, while the Canberra light rail is scheduled to open in April 2019.[54] This will be the first time that there have been trams in Canberra, even though Walter Burley Griffin's 1914-1920 plans for the capital then in the planning stage did propose a Canberra tram system.[55]
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+
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+ In Japan, the Kyoto Electric railroad was the first tram system, starting operation in 1895.[56] By 1932, the network had grown to 82 railway companies in 65 cities, with a total network length of 1,479 km (919 mi).[57] By the 1960s the tram had generally died out in Japan.[58][59]
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+
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+ Two rare but significant alternatives were conduit current collection, which was widely used in London, Washington, D.C. and New York City, and the surface contact collection method, used in Wolverhampton (the Lorain system), Torquay and Hastings in the UK (the Dolter stud system), and currently in Bordeaux, France (the ground-level power supply system).[citation needed]
90
+
91
+ The convenience and economy of electricity resulted in its rapid adoption once the technical problems of production and transmission of electricity were solved. Electric trams largely replaced animal power and other forms of motive power including cable and steam, in the late 19th and early 20th centuries.[citation needed]
92
+
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+ There is one particular hazard associated with trams powered from a trolley pole off an overhead line. Since the tram relies on contact with the rails for the current return path, a problem arises if the tram is derailed or (more usually) if it halts on a section of track that has been particularly heavily sanded by a previous tram, and the tram loses electrical contact with the rails. In this event, the underframe of the tram, by virtue of a circuit path through ancillary loads (such as interior lighting), is live at the full supply voltage, typically 600 volts DC. In British terminology, such a tram was said to be ‘grounded’—not to be confused with the US English use of the term, which means the exact opposite. Any person stepping off the tram completed the earth return circuit and could receive a nasty electric shock. In such an event, the driver was required to jump off the tram (avoiding simultaneous contact with the tram and the ground) and pull down the trolley pole, before allowing passengers off the tram. Unless derailed, the tram could usually be recovered by running water down the running rails from a point higher than the tram, the water providing a conducting bridge between the tram and the rails.[citation needed]
94
+
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+ In the 2000s, two companies introduced catenary-free designs. Alstom's Citadis line uses a third rail, and Bombardier's PRIMOVE LRV is charged by contactless induction plates embedded in the trackway.[60]
96
+
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+ In some places, other forms of power were used to power the tram.
98
+
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+ As early as 1834, Thomas Davenport, a Vermont blacksmith, had invented a battery-powered electric motor which he later patented. The following year he used it to operate a small model electric car on a short section of track four feet in diameter.[61][62]
100
+
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+ Attempts to use batteries as a source of electricity were made from the 1880s and 1890s, with unsuccessful trials conducted in among other places Bendigo and Adelaide in Australia, and for about 14 years as The Hague accutram of HTM in the Netherlands. The first trams in Bendigo, Australia, in 1892, were battery-powered but within as little as three months they were replaced with horse-drawn trams. In New York City some minor lines also used storage batteries. Then, comparatively recently, during the 1950s, a longer battery-operated tramway line ran from Milan to Bergamo. In China there is a Nanjing battery Tram line and has been running since 2014.[63] More recently in 2019, the West Midlands Metro in Birmingham, England has adopted battery powered trams on sections through the city centre close to Grade I listed Birmingham Town Hall.
102
+
103
+ Paris and Berne (Switzerland) [64][circular reference] operated trams that were powered by compressed air using the Mekarski system.
104
+
105
+ The Convict Tramway [65] was hauled by human power in the form of convicts from the Port Arthur convict settlement.[66] and was created to replace the hazardous sea voyage from Hobart to Port Arthur, Tasmania.[67][65] Charles O'Hara Booth oversaw the construction of the tramway.[68]
106
+
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+ It opened in 1836 and ran for 8 km (5 miles) from Oakwood to Taranna.[69] By most definitions, the tramway was the first passenger-carrying railway/tramway in Australia.[67] An unconfirmed report says that it continued to Eaglehawk Neck and, if this was so, the length of the tramway would have been more than doubled. The tramway carried passengers and freight, and ran on wooden rails. The gauge is unknown. The date of closure is unknown, but it was certainly prior to 1877.[70]
108
+
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+ In March 2015, China South Rail Corporation (CSR) demonstrated the world's first hydrogen fuel cell vehicle tramcar at an assembly facility in Qingdao. The chief engineer of the CSR subsidiary CSR Sifang Co Ltd., Liang Jianying, said that the company is studying how to reduce the running costs of the tram.[71][72]
110
+
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+ The Trieste–Opicina tramway in Trieste operates a hybrid funicular tramway system. Conventional electric trams are operated in street running and on reserved track for most of their route. However, on one steep segment of track, they are assisted by cable tractors, which push the trams uphill and act as brakes for the downhill run. For safety, the cable tractors are always deployed on the downhill side of the tram vehicle.
112
+
113
+ Similar systems were used elsewhere in the past, notably on the Queen Anne Counterbalance in Seattle and the Darling Street wharf line in Sydney.
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+
115
+ Hastings and some other tramways, for example Stockholms Spårvägar in Sweden and some lines in Karachi, used petrol trams. Galveston Island Trolley in Texas operated diesel trams due to the city's hurricane-prone location, which would result in frequent damage to an electrical supply system.
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+
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+ Although Portland, Victoria promotes its tourist tram[73] as being a cable car it actually operates using a hidden diesel motor. The tram, which runs on a circular route around the town of Portland, uses dummies and salons formerly used on the extensive Melbourne cable tramway system and now beautifully restored.
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+
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+ In the mid-20th century many tram systems were disbanded, replaced by buses, automobiles or rapid transit. The General Motors streetcar conspiracy was a case study of the decline of trams in the United States. In the 21st century, trams have been re-introduced in cities where they had been closed down for decades (such as Tramlink in London), or kept in heritage use (such as Spårväg City in Stockholm). Vehicle fabricates from the 1990s and onwards (such as the Bombardier Flexity series and Alstom Citadis) are usually low-floor trams with features such as articulation and regenerative braking.
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+ Trams have been used for two main purposes: for carrying passengers and for carrying cargo. There are several types of passenger tram:
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+ There are two main types of tramways, the classic tramway built in the early 20th century with the tram system operating in mixed traffic, and the later type which is most often associated with the tram system having its own right of way. Tram systems that have their own right of way are often called light rail but this does not always hold true. Though these two systems differ in their operation, their equipment is much the same.
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+ Trams were traditionally operated with separate levers for applying power and brakes. More modern vehicles use a locomotive-style controller which incorporate a dead man's switch. The success of the PCC streetcar had also seen trams use automobile-style foot controls allowing hands-free operation, particularly when the driver was responsible for fare collection.
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+
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+ Electric trams use various devices to collect power from overhead lines. The most common device found today is the pantograph, while some older systems use trolley poles or bow collectors. Ground-level power supply has become a recent innovation. Another new technology uses supercapacitors; when an insulator at a track switch cuts off power from the tram for a short distance along the line, the tram can use energy stored in a large capacitor to drive the tram past the gap in the power feed.[74] A rather obsolete system for power supply is conduit current collection.
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+
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+ The old tram systems in London, Manhattan (New York City), and Washington, D.C., used live rails, like those on third-rail electrified railways, but in a conduit underneath the road, from which they drew power through a plough. It was called Conduit current collection. Washington's was the last of these to close, in 1962. Today, no commercial tramway uses this system. More recently, a modern equivalent to these systems has been developed which allows for the safe installation of a third rail on city streets, which is known as surface current collection or ground-level power supply; the main example of this is the new tramway in Bordeaux.
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+
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+ A ground-level power supply system also known as Surface current collection or Alimentation par le sol (APS) is an updated version of the original stud type system. APS uses a third rail placed between the running rails, divided electrically into eight-metre powered segments with three metre neutral sections between. Each tram has two power collection skates, next to which are antennas that send radio signals to energize the power rail segments as the tram passes over them.
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+
133
+ Older systems required mechanical switching systems which were susceptible to environmental problems. At any one time no more than two consecutive segments under the tram should actually be live. Wireless and solid state switching remove the mechanical problem.
134
+
135
+ Alstom developed the system primarily to avoid intrusive power supply cables in the sensitive area of the old city of old Bordeaux.[75]
136
+
137
+ Route patterns vary greatly among the world's tram systems, leading to different network topologies.
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+
139
+ The resulting route patterns are very different. Some have a rational structure, covering their catchment area as efficiently as possible, with new suburbs being planned with tramlines integral to their layout – such is the case in Amsterdam. Bordeaux and Montpellier have built comprehensive networks, based on radial routes with numerous interconnections, within the last two decades. Some systems serve only parts of their cities, with Berlin being the prime example, owing to the fact that trams survived the city's political division only in the Eastern part. Other systems have ended up with a rather random route map, for instance when some previous operating companies have ceased operation (as with the tramways vicinaux/buurtspoorwegen in Brussels) or where isolated outlying lines have been preserved (as on the eastern fringe of Berlin). In Rome, the remnant of the system comprises 3 isolated radial routes, not connecting in the ancient city centre, but linked by a ring route. Some apparently anomalous lines continue in operation where a new line would not on rational grounds be built, because it is much more costly to build a new line than continue operating an existing one.
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+
141
+ In some places, the opportunity is taken when roads are being repaved to lay tramlines (though without erecting overhead cables) even though no service is immediately planned: such is the case in Leipzigerstraße in Berlin, the Haarlemmer Houttuinen in Amsterdam, and Botermarkt in Ghent.
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+
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+ Tram systems operate across national borders in Basel (from Switzerland into France and Germany) and Strasbourg (From France into Germany). It is planned to open a line linking Hasselt (Belgium) with Maastricht (Netherlands) in 2021.
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+
145
+ Tramway track can have different rail profiles to accommodate the various operating environments of the vehicle. They may be embedded into concrete for street-running operation, or use standard ballasted track with railroad ties on high-speed sections. A more ecological solution is to embed tracks into grass turf.
146
+
147
+ Tramway tracks use a grooved rail with a groove designed for tramway or railway track in pavement or grassed surfaces (grassed track or track in a lawn). The rail has the railhead on one side and the guard on the other. The guard provides accommodation for the flange. The guard carries no weight, but may act as a checkrail. Grooved rail was invented in 1852 by Alphonse Loubat, a French inventor who developed improvements in tram and rail equipment, and helped develop tram lines in New York City and Paris. The invention of grooved rail enabled tramways to be laid without causing a nuisance to other road users, except unsuspecting cyclists, who could get their wheels caught in the groove. The grooves may become filled with gravel and dirt (particularly if infrequently used or after a period of idleness) and need clearing from time to time, this being done by a "scrubber" tram. Failure to clear the grooves can lead to a bumpy ride for the passengers, damage to either wheel or rail and possibly derailing.
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+ In narrow situations double-track tram lines sometimes reduce to single track, or, to avoid switches, have the tracks interlaced, e.g. in the Leidsestraat in Amsterdam on three short stretches (see map detail); this is known as interlaced or gauntlet track. There is a UK example of interlaced track on the Tramlink, just west of Mitcham Station, where the formation is narrowed by an old landslip causing an obstruction. (See photo in Tramlink entry).
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+
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+ Historically, the track gauge has had considerable variations, with narrow gauge common in many early systems. However, most light rail systems are now standard gauge. An important advantage of standard gauge is that standard railway maintenance equipment can be used on it, rather than custom-built machinery. Using standard gauge also allows light rail vehicles to be delivered and relocated conveniently using freight railways and locomotives.
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+
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+ Another factor favoring standard gauge is that low-floor vehicles are becoming popular, and there is generally insufficient space for wheelchairs to move between the wheels in a narrow gauge layout. Standard gauge also enables – at least in theory – a larger choice of manufacturers and thus lower procurement costs for new vehicles. However, other factors such as electrification or loading gauge for which there is more variation may require costly custom built units regardless.
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+ Tram stops may be similar to bus stops in design and use, particularly in street-running sections, where in some cases other vehicles are legally required to stop clear of the tram doors. Some stops may resemble to railway platforms, particularly in private right-of-way sections and where trams are boarded at standard railway platform height, as opposed to using steps at the doorway or low-floor trams.
156
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+ Approximately 5,000 new trams are manufactured each year. As of February 2017, 4,478 new trams were on order from their makers, with options being open for a further 1,092.[76]
158
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+ The main manufacturers are:
160
+
161
+ Trams are in a period of growth, with about 800 tram systems operating around the world, 10 or so new systems being opened each year, and many being gradually extended.[87] Some of these systems date from the late 19th or early 20th centuries. In the past 20 years their numbers have been augmented by modern tramway or light rail systems in cities that had discarded this form of transport. There have also been some new tram systems in cities that never previously had them.
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+ Tramways with tramcars (British English) or street railways with streetcars (North American English) were common throughout the industrialised world in the late 19th and early 20th centuries but they had disappeared from most British, Canadian, French and US cities by the mid-20th century.[88]
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+ By contrast, trams in parts of continental Europe continued to be used by many cities, although there were contractions in some countries, including the Netherlands.[89]
166
+
167
+ Since 1980 trams have returned to favour in many places, partly because their tendency to dominate the roadway, formerly seen as a disadvantage, is now considered to be a merit since it raises the visibility of public transport (encouraging car users to change their mode of travel), and enables streets to be reconfigured to give more space to pedestrians, making cites more pleasant places to live. New systems have been built in the United States, United Kingdom, Ireland, Italy, France, Australia and many other countries.
168
+
169
+ In Milan, Italy, the old "Ventotto" trams are considered by its inhabitants a "symbol" of the city. The same can be said of trams in Melbourne in general, but particularly the iconic W class. The Toronto streetcar system had similarly become an iconic symbol of the city, operating the largest network in the Americas as well as the only large-scale tram system in Canada (not including light rail systems, or heritage lines).[90][91]
170
+
171
+ The largest tram ((classic tram, streetcar, straßenbahn) and fast tram (light rail, stadtbahn)) networks in the world by route length (as of 2016)[92] are:
172
+
173
+ Other large transit networks that operate streetcar and light rail systems include:
174
+
175
+ This list is not exhaustive.
176
+
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+
178
+
179
+ Historically, the Paris Tram System was, at its peak, the world's largest system, with 1,111 km (690 mi) of track in 1925[citation needed] (according to other sources, ca. 640 km (400 mi) of route length in 1930). However it was completely closed in 1938.[133] The next largest system appears to have been 857 km (533 mi), in Buenos Aires before 19 February 1963. The third largest was Chicago, with over 850 km (530 mi) of track,[134] but it was all converted to trolleybus and bus services by 21 June 1958. Before its decline, the BVG in Berlin operated a very large network with 634 km (394 mi) of route. Before its system started to be converted to trolleybus (and later bus) services in the 1930s (last tramway closed 6 July 1952), the first-generation London network had 555 km (345 mi) of route in 1931.[135] In 1958 trams in Rio de Jainero were employed on (433 km; 269 mi) of track. The final line, the Santa teresa route was closed in 1968.[136] During a period in the 1980s, the world's largest tram system was in Leningrad (now known as St. Petersburg) with 350 km (220 mi), USSR, and was included as such in the Guinness World Records;[citation needed] however Saint Petersburg's tram system has declined in size since the fall of the Soviet Union. Vienna in 1960 had 340 km (211 mi), before the expansion of bus services and the opening of a subway (1976). Substituting subway services for tram routes continues. 320 km (199 mi) was in Minneapolis-Saint Paul in 1947: There streetcars ended 31 October 1953 in Minneapolis and 19 June 1954 in St. Paul.[137] The Sydney tram network, before it was closed on 25 February 1961, had 291 km (181 mi) of route, and was thus the largest in Australia. As from 1961, the Melbourne system (currently recognised as the world's largest) took over Sydney's title as the largest network in Australia.
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+
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+ In many European cities, much tramway infrastructure was lost in the mid-20th century, though not always on the same scale as in other parts of the world such as North America. Most of Central and Eastern Europe retained the majority of its tramway systems and it is here that the largest and busiest tram systems in the world are found.
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+ Whereas most systems and vehicles in the tram sector are found in Central and Eastern Europe, in the 1960s and 1970s, tram systems were shut down in many places in Western Europe, however urban transportation has been experiencing a sustained long running revival since the 1990s. Many European cities are rehabilitating, upgrading, expanding and reconstructing their old tramway lines and building new tramway lines.[139]
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+ In North America, these vehicles are called "streetcars" (or "trolleys"); the term tram is more likely to be understood as an aerial tramway or a people-mover. Streetcar systems were developed in late 19th to early 20th centuries in a number of cities throughout North America. However, most North American cities saw its streetcar lines removed in the mid-20th century for a variety of financial, technological and social reasons. Exceptions included Boston,[140] Cleveland, Mexico City, New Orleans, Newark, Philadelphia, Pittsburgh, San Francisco, and Toronto.
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+ Toronto currently operates the largest streetcar system in the Americas in terms of track length and ridership. Operated by the Toronto Transit Commission, the streetcar system is the only large-scale streetcar system existing in Canada, excluding heritage streetcar, or light rail systems that are operated in other Canadian municipalities. The streetcar system was established in 1861, and used a variety of vehicles in its history, including horse-drawn streetcars, Peter Witt streetcars, the PCC streetcar, and the Canadian Light Rail Vehicle and its articulated counterpart, the Articulated Light Rail Vehicle. Since December 29, 2019,[141] the system exclusively uses the Flexity Outlook made by Bombardier Transportation.[142][143][144][145]
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+
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+ Streetcars once existed in the Canadian cities of Calgary, Edmonton, Halifax, Hamilton, Kingston, Kitchener, London, Montreal, Ottawa, Peterborough, Quebec City, Regina, Saskatoon, Windsor, and Vancouver. However, Canadian cities excluding Toronto, removed their streetcar systems in the mid-20th century. In the late 1970s and early 1980s, light rail systems were introduced in Calgary and Edmonton; with another light rail system established in Ottawa in 2001. There is now something of a renaissance for light railways in mid-sized cities with Waterloo, Ontario the first to come on line and construction underway in Mississauga, Ontario. In the late 20th century, several Canadian locales restored portions of their defunct streetcar lines, operating them as a heritage feature for tourists. Heritage streetcar lines in Canada include the High Level Bridge Streetcar in Edmonton, the Nelson Electric Tramway in Nelson, and the Whitehorse Waterfront Trolley in Whitehorse.
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+ Pittsburgh had kept most of its streetcar system serving the city and many suburbs, making it the longest-lasting large-network streetcar system in the United States.[citation needed] However, most of Pittsburgh's surviving streetcar lines were converted to light rail in the 1980s. San Francisco's Muni Metro system is the largest surviving streetcar system in the United States, and has even revived previously closed streetcar lines such as the F Market & Wharves heritage streetcar line. In the late 20th century, several cities installed modern light rail systems, in part along the same corridors as their old streetcars systems, the first of these being the San Diego Trolley in San Diego in 1981.
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+ In the 1980s, some cities in the United States brought back streetcars lines, including Memphis, Tampa, and Little Rock; However, these streetcar systems were designed as heritage streetcar lines, and used vintage or replica-vintage vehicles. The first "second-generation streetcar systems" in North America was opened in Portland in 2001.[146] The "second-generation streetcar system," utilizes modern vehicles – vehicles that feature low-floor streetcars. These newer streetcar systems were built in several American cities in the early 21st century including Atlanta, Charlotte, Cincinnati, Dallas, Detroit, Kansas City, Milwaukee, Oklahoma City, Seattle, Tucson, and Washington, D.C..
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+ Model trams are popular in HO scale (1:87) and O scale (1:48 in the US and generally 1:43,5 and 1:45 in Europe and Asia). They are typically powered and will accept plastic figures inside. Common manufacturers are Roco and Lima, with many custom models being made as well. The German firm Hödl[165] and the Austrian Halling[166] specialise in 1:87 scale.[167]
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+
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+ In the US, Bachmann Industries is a mass supplier of HO streetcars and kits. Bowser Manufacturing has produced white metal models for over 50 years.[168] There are many boutique vendors offering limited run epoxy and wood models. At the high end are highly detailed brass models which are usually imported from Japan or Korea and can cost in excess of $500. Many of these run on 16.5 mm (0.65 in) gauge track, which is correct for the representation of 4 ft 8 1⁄2 in (1,435 mm) (standard gauge) in HO scale as in US and Japan, but incorrect in 4 mm (1:76.2) scale, as it represents 4 ft 8 1⁄2 in (1,435 mm). This scale/gauge hybrid is called OO scale.
198
+ O scale trams are also very popular among tram modellers because the increased size allows for more detail and easier crafting of overhead wiring. In the US these models are usually purchased in epoxy or wood kits and some as brass models. The Saint Petersburg Tram Company[169] produces highly detailed polyurethane non-powered O Scale models from around the world which can easily be powered by trucks from vendors like Q-Car.[170]
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+ In the US, one of the best resources for model tram enthusiasts is the East Penn Traction Club of Philadelphia [171] and Trolleyville a website of the Southern California Traction Club.[172]
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+ It is thought that the first example of a working model tramcar in the UK built by an amateur for fun was in 1929, when Frank E. Wilson created a replica of London County Council Tramways E class car 444 in 1:16 scale, which he demonstrated at an early Model Engineer Exhibition. Another of his models was London E/1 1800, which was the only tramway exhibit in the Faraday Memorial Exhibition of 1931. Together with likeminded friends, Frank Wilson went on to found the Tramway & Light Railway Society[173] in 1938, establishing tramway modelling as a hobby.
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+ Mr. Nathan was a passenger by No. 2 tramway car [...] [he] alighted from the car at the southern end, but before he got clear of the rails the car moved onwards [...] he was thus whirled round by the sudden motion of the carriage and his body was brought under the front wheel.
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1
+
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+
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+ Rail transport (also known as train transport) is a means of transferring passengers and goods on wheeled vehicles running on rails, which are located on tracks. In contrast to road transport, where vehicles run on a prepared flat surface, rail vehicles (rolling stock) are directionally guided by the tracks on which they run. Tracks usually consist of steel rails, installed on ties (sleepers) set in ballast, on which the rolling stock, usually fitted with metal wheels, moves. Other variations are also possible, such as slab track. This is where the rails are fastened to a concrete foundation resting on a prepared subsurface.
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+ Rolling stock in a rail transport system generally encounters lower frictional resistance than rubber-tired road vehicles, so passenger and freight cars (carriages and wagons) can be coupled into longer trains. The operation is carried out by a railway company, providing transport between train stations or freight customer facilities. Power is provided by locomotives which either draw electric power from a railway electrification system or produce their own power, usually by diesel engines or, historically, steam engines. Most tracks are accompanied by a signalling system. Railways are a safe land transport system when compared to other forms of transport.[Nb 1] Railway transport is capable of high levels of passenger and cargo utilization and energy efficiency, but is often less flexible and more capital-intensive than road transport, when lower traffic levels are considered.
6
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+ The oldest known, man/animal-hauled railways date back to the 6th century BC in Corinth, Greece. Rail transport then commenced in mid 16th century in Germany in the form of horse-powered funiculars and wagonways. Modern rail transport commenced with the British development of the steam locomotives in the early 19th century. Thus the railway system in Great Britain is the oldest in the world. Built by George Stephenson and his son Robert's company Robert Stephenson and Company, the Locomotion No. 1 is the first steam locomotive to carry passengers on a public rail line, the Stockton and Darlington Railway in 1825. George Stephenson also built the first public inter-city railway line in the world to use only the steam locomotives all the time, the Liverpool and Manchester Railway which opened in 1830. With steam engines, one could construct mainline railways, which were a key component of the Industrial Revolution. Also, railways reduced the costs of shipping, and allowed for fewer lost goods, compared with water transport, which faced occasional sinking of ships. The change from canals to railways allowed for "national markets" in which prices varied very little from city to city. The spread of the railway network and the use of railway timetables, led to the standardisation of time (railway time) in Britain based on Greenwich Mean Time. Prior to this, major towns and cities varied their local time relative to GMT. The invention and development of the railway in the United Kingdom was one of the most important technological inventions of the 19th century. The world's first underground railway, the Metropolitan Railway (part of the London Underground), opened in 1863.
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+ In the 1880s, electrified trains were introduced, leading to electrification of tramways and rapid transit systems. Starting during the 1940s, the non-electrified railways in most countries had their steam locomotives replaced by diesel-electric locomotives, with the process being almost complete by the 2000s. During the 1960s, electrified high-speed railway systems were introduced in Japan and later in some other countries. Many countries are in the process of replacing diesel locomotives with electric locomotives, mainly due to environmental concerns, a notable example being Switzerland, which has completely electrified its network. Other forms of guided ground transport outside the traditional railway definitions, such as monorail or maglev, have been tried but have seen limited use.
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+ Following a decline after World War II due to competition from cars and airplanes, rail transport has had a revival in recent decades due to road congestion and rising fuel prices, as well as governments investing in rail as a means of reducing CO2 emissions in the context of concerns about global warming.
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+ The history of rail transport began in the 6th century BC in Ancient Greece. It can be divided up into several discrete periods defined by the principal means of track material and motive power used.
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+ Evidence indicates that there was 6 to 8.5 km long Diolkos paved trackway, which transported boats across the Isthmus of Corinth in Greece from around 600 BC.[1][2][3][4][5] Wheeled vehicles pulled by men and animals ran in grooves in limestone, which provided the track element, preventing the wagons from leaving the intended route. The Diolkos was in use for over 650 years, until at least the 1st century AD.[5] Paved trackways were also later built in Roman Egypt.[6]
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+ In 1515, Cardinal Matthäus Lang wrote a description of the Reisszug, a funicular railway at the Hohensalzburg Fortress in Austria. The line originally used wooden rails and a hemp haulage rope and was operated by human or animal power, through a treadwheel.[7] The line still exists and is operational, although in updated form and is possibly the oldest operational railway.[8]
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+ Wagonways (or tramways) using wooden rails, hauled by horses, started appearing in the 1550s to facilitate the transport of ore tubs to and from mines, and soon became popular in Europe. Such an operation was illustrated in Germany in 1556 by Georgius Agricola in his work De re metallica.[9] This line used "Hund" carts with unflanged wheels running on wooden planks and a vertical pin on the truck fitting into the gap between the planks to keep it going the right way. The miners called the wagons Hunde ("dogs") from the noise they made on the tracks.[10]
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+ There are many references to their use in central Europe in the 16th century.[11] Such a transport system was later used by German miners at Caldbeck, Cumbria, England, perhaps from the 1560s.[12] A wagonway was built at Prescot, near Liverpool, sometime around 1600, possibly as early as 1594. Owned by Philip Layton, the line carried coal from a pit near Prescot Hall to a terminus about half a mile away.[13] A funicular railway was also made at Broseley in Shropshire some time before 1604. This carried coal for James Clifford from his mines down to the river Severn to be loaded onto barges and carried to riverside towns.[14] The Wollaton Wagonway, completed in 1604 by Huntingdon Beaumont, has sometimes erroneously been cited as the earliest British railway. It ran from Strelley to Wollaton near Nottingham.[15]
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+
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+ The Middleton Railway in Leeds, which was built in 1758, later became the world's oldest operational railway (other than funiculars), albeit now in an upgraded form. In 1764, the first railway in the Americas was built in Lewiston, New York.[16]
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+
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+ In the late 1760s, the Coalbrookdale Company began to fix plates of cast iron to the upper surface of the wooden rails. This allowed a variation of gauge to be used. At first only balloon loops could be used for turning, but later, movable points were taken into use that allowed for switching.[17]
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+ A system was introduced in which unflanged wheels ran on L-shaped metal plates – these became known as plateways. John Curr, a Sheffield colliery manager, invented this flanged rail in 1787, though the exact date of this is disputed. The plate rail was taken up by Benjamin Outram for wagonways serving his canals, manufacturing them at his Butterley ironworks. In 1803, William Jessop opened the Surrey Iron Railway, a double track plateway, erroneously sometimes cited as world's first public railway, in south London.[18]
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+ Meanwhile, William Jessop had earlier used a form of all-iron edge rail and flanged wheels successfully for an extension to the Charnwood Forest Canal at Nanpantan, Loughborough, Leicestershire in 1789. In 1790, Jessop and his partner Outram began to manufacture edge-rails. Jessop became a partner in the Butterley Company in 1790. The first public edgeway (thus also first public railway) built was Lake Lock Rail Road in 1796. Although the primary purpose of the line was to carry coal, it also carried passengers.
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+
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+ These two systems of constructing iron railways, the "L" plate-rail and the smooth edge-rail, continued to exist side by side until well into the early 19th century. The flanged wheel and edge-rail eventually proved its superiority and became the standard for railways.
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+ Cast iron used in rails proved unsatisfactory because it was brittle and broke under heavy loads. The wrought iron invented by John Birkinshaw in 1820 replaced cast iron. Wrought iron (usually simply referred to as "iron") was a ductile material that could undergo considerable deformation before breaking, making it more suitable for iron rails. But iron was expensive to produce until Henry Cort patented the puddling process in 1784. In 1783 Cort also patented the rolling process, which was 15 times faster at consolidating and shaping iron than hammering.[19] These processes greatly lowered the cost of producing iron and rails. The next important development in iron production was hot blast developed by James Beaumont Neilson (patented 1828), which considerably reduced the amount of coke (fuel) or charcoal needed to produce pig iron.[20] Wrought iron was a soft material that contained slag or dross. The softness and dross tended to make iron rails distort and delaminate and they lasted less than 10 years. Sometimes they lasted as little as one year under high traffic. All these developments in the production of iron eventually led to replacement of composite wood/iron rails with superior all iron rails.
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+ The introduction of the Bessemer process, enabling steel to be made inexpensively, led to the era of great expansion of railways that began in the late 1860s. Steel rails lasted several times longer than iron.[21][22][23] Steel rails made heavier locomotives possible, allowing for longer trains and improving the productivity of railroads.[24] The Bessemer process introduced nitrogen into the steel, which caused the steel to become brittle with age. The open hearth furnace began to replace the Bessemer process near the end of the 19th century, improving the quality of steel and further reducing costs. Thus steel completely replaced the use of iron in rails, becoming standard for all railways.
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+ The first passenger horsecar or tram, Swansea and Mumbles Railway was opened between Swansea and Mumbles in Wales in 1807.[25] Horses remained the preferable mode for tram transport even after the arrival of steam engines until the end of the 19th century, because they were cleaner compared to steam driven trams which caused smoke in city streets.
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+ In 1784 James Watt, a Scottish inventor and mechanical engineer, patented a design for a steam locomotive. Watt had improved the steam engine of Thomas Newcomen, hitherto used to pump water out of mines, and developed a reciprocating engine in 1769 capable of powering a wheel. This was a large stationary engine, powering cotton mills and a variety of machinery; the state of boiler technology necessitated the use of low pressure steam acting upon a vacuum in the cylinder, which required a separate condenser and an air pump. Nevertheless, as the construction of boilers improved, Watt investigated the use of high-pressure steam acting directly upon a piston, raising the possibility of a smaller engine that might be used to power a vehicle. Following his patent, Watt's employee William Murdoch produced a working model of a self-propelled steam carriage in that year.[26]
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+ The first full-scale working railway steam locomotive was built in the United Kingdom in 1804 by Richard Trevithick, a British engineer born in Cornwall. This used high-pressure steam to drive the engine by one power stroke. The transmission system employed a large flywheel to even out the action of the piston rod. On 21 February 1804, the world's first steam-powered railway journey took place when Trevithick's unnamed steam locomotive hauled a train along the tramway of the Penydarren ironworks, near Merthyr Tydfil in South Wales.[27][28] Trevithick later demonstrated a locomotive operating upon a piece of circular rail track in Bloomsbury, London, the Catch Me Who Can, but never got beyond the experimental stage with railway locomotives, not least because his engines were too heavy for the cast-iron plateway track then in use.[29]
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+ The first commercially successful steam locomotive was Matthew Murray's rack locomotive Salamanca built for the Middleton Railway in Leeds in 1812. This twin-cylinder locomotive was light enough to not break the edge-rails track and solved the problem of adhesion by a cog-wheel using teeth cast on the side of one of the rails. Thus it was also the first rack railway.
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+
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+ This was followed in 1813 by the locomotive Puffing Billy built by Christopher Blackett and William Hedley for the Wylam Colliery Railway, the first successful locomotive running by adhesion only. This was accomplished by the distribution of weight between a number of wheels. Puffing Billy is now on display in the Science Museum in London, making it the oldest locomotive in existence.[30]
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+
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+ In 1814 George Stephenson, inspired by the early locomotives of Trevithick, Murray and Hedley, persuaded the manager of the Killingworth colliery where he worked to allow him to build a steam-powered machine. Stephenson played a pivotal role in the development and widespread adoption of the steam locomotive. His designs considerably improved on the work of the earlier pioneers. He built the locomotive Blücher, also a successful flanged-wheel adhesion locomotive. In 1825 he built the locomotive Locomotion for the Stockton and Darlington Railway in the north east of England, which became the first public steam railway in the world in 1825, although it used both horse power and steam power on different runs. In 1829, he built the locomotive Rocket, which entered in and won the Rainhill Trials. This success led to Stephenson establishing his company as the pre-eminent builder of steam locomotives for railways in Great Britain and Ireland, the United States, and much of Europe.[31]:24–30 The first public railway which used only steam locomotives, all the time, was Liverpool and Manchester Railway, built in 1830.
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+ Steam power continued to be the dominant power system in railways around the world for more than a century.
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+ The first known electric locomotive was built in 1837 by chemist Robert Davidson of Aberdeen in Scotland, and it was powered by galvanic cells (batteries). Thus it was also the earliest battery electric locomotive. Davidson later built a larger locomotive named Galvani, exhibited at the Royal Scottish Society of Arts Exhibition in 1841. The seven-ton vehicle had two direct-drive reluctance motors, with fixed electromagnets acting on iron bars attached to a wooden cylinder on each axle, and simple commutators. It hauled a load of six tons at four miles per hour (6 kilometers per hour) for a distance of one and a half miles (2.4 kilometres). It was tested on the Edinburgh and Glasgow Railway in September of the following year, but the limited power from batteries prevented its general use. It was destroyed by railway workers, who saw it as a threat to their job security.[32][33][34]
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+
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+ Werner von Siemens demonstrated an electric railway in 1879 in Berlin. The world's first electric tram line, Gross-Lichterfelde Tramway, opened in Lichterfelde near Berlin, Germany, in 1881. It was built by Siemens. The tram ran on 180 Volt DC, which was supplied by running rails. In 1891 the track was equipped with an overhead wire and the line was extended to Berlin-Lichterfelde West station. The Volk's Electric Railway opened in 1883 in Brighton, England. The railway is still operational, thus making it the oldest operational electric railway in the world. Also in 1883, Mödling and Hinterbrühl Tram opened near Vienna in Austria. It was the first tram line in the world in regular service powered from an overhead line. Five years later, in the U.S. electric trolleys were pioneered in 1888 on the Richmond Union Passenger Railway, using equipment designed by Frank J. Sprague.[35]
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+
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+ The first use of electrification on a main line was on a four-mile section of the Baltimore Belt Line of the Baltimore and Ohio Railroad (B&O) in 1895 connecting the main portion of the B&O to the new line to New York through a series of tunnels around the edges of Baltimore's downtown. Electricity quickly became the power supply of choice for subways, abetted by the Sprague's invention of multiple-unit train control in 1897. By the early 1900s most street railways were electrified.
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+ The London Underground, the world's oldest underground railway, opened in 1863, and it began operating electric services using a fourth rail system in 1890 on the City and South London Railway, now part of the London Underground Northern line. This was the first major railway to use electric traction. The world's first deep-level electric railway, it runs from the City of London, under the River Thames, to Stockwell in south London.[36]
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+ The first practical AC electric locomotive was designed by Charles Brown, then working for Oerlikon, Zürich. In 1891, Brown had demonstrated long-distance power transmission, using three-phase AC, between a hydro-electric plant at Lauffen am Neckar and Frankfurt am Main West, a distance of 280 km. Using experience he had gained while working for Jean Heilmann on steam-electric locomotive designs, Brown observed that three-phase motors had a higher power-to-weight ratio than DC motors and, because of the absence of a commutator, were simpler to manufacture and maintain.[37] However, they were much larger than the DC motors of the time and could not be mounted in underfloor bogies: they could only be carried within locomotive bodies.[38]
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+ In 1894, Hungarian engineer Kálmán Kandó developed a new type 3-phase asynchronous electric drive motors and generators for electric locomotives. Kandó's early 1894 designs were first applied in a short three-phase AC tramway in Evian-les-Bains (France), which was constructed between 1896 and 1898.[39][40][41][42][43]
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+
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+ In 1896, Oerlikon installed the first commercial example of the system on the Lugano Tramway. Each 30-tonne locomotive had two 110 kW (150 hp) motors run by three-phase 750 V 40 Hz fed from double overhead lines. Three-phase motors run at constant speed and provide regenerative braking, and are well suited to steeply graded routes, and the first main-line three-phase locomotives were supplied by Brown (by then in partnership with Walter Boveri) in 1899 on the 40 km Burgdorf–Thun line, Switzerland.
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+ Italian railways were the first in the world to introduce electric traction for the entire length of a main line rather than a short section. The 106 km Valtellina line was opened on 4 September 1902, designed by Kandó and a team from the Ganz works.[44][45] The electrical system was three-phase at 3 kV 15 Hz. In 1918,[46] Kandó invented and developed the rotary phase converter, enabling electric locomotives to use three-phase motors whilst supplied via a single overhead wire, carrying the simple industrial frequency (50 Hz) single phase AC of the high voltage national networks.[45]
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+ An important contribution to the wider adoption of AC traction came from SNCF of France after World War II. The company conducted trials at AC 50 Hz, and established it as a standard. Following SNCF's successful trials, 50 Hz, now also called industrial frequency was adopted as standard for main-lines across the world.[47]
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+ Earliest recorded examples of an internal combustion engine for railway use included a prototype designed by William Dent Priestman, which was examined by Sir William Thomson in 1888 who described it as a "[Priestman oil engine] mounted upon a truck which is worked on a temporary line of rails to show the adaptation of a petroleum engine for locomotive purposes.".[48][49] In 1894, a 20 hp (15 kW) two axle machine built by Priestman Brothers was used on the Hull Docks.[50]
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+
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+ In 1906, Rudolf Diesel, Adolf Klose and the steam and diesel engine manufacturer Gebrüder Sulzer founded Diesel-Sulzer-Klose GmbH to manufacture diesel-powered locomotives. Sulzer had been manufacturing diesel engines since 1898. The Prussian State Railways ordered a diesel locomotive from the company in 1909. The world's first diesel-powered locomotive was operated in the summer of 1912 on the Winterthur–Romanshorn railway in Switzerland, but was not a commercial success.[51] The locomotive weight was 95 tonnes and the power was 883 kW with a maximum speed of 100 km/h.[52] Small numbers of prototype diesel locomotives were produced in a number of countries through the mid-1920s.
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+ A significant breakthrough occurred in 1914, when Hermann Lemp, a General Electric electrical engineer, developed and patented a reliable direct current electrical control system (subsequent improvements were also patented by Lemp).[53] Lemp's design used a single lever to control both engine and generator in a coordinated fashion, and was the prototype for all diesel–electric locomotive control systems. In 1914, world's first functional diesel–electric railcars were produced for the Königlich-Sächsische Staatseisenbahnen (Royal Saxon State Railways) by Waggonfabrik Rastatt with electric equipment from Brown, Boveri & Cie and diesel engines from Swiss Sulzer AG. They were classified as DET 1 and DET 2 (de.wiki). The first regular use of diesel–electric locomotives was in switching (shunter) applications. General Electric produced several small switching locomotives in the 1930s (the famous "44-tonner" switcher was introduced in 1940) Westinghouse Electric and Baldwin collaborated to build switching locomotives starting in 1929.
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+ In 1929, the Canadian National Railways became the first North American railway to use diesels in mainline service with two units, 9000 and 9001, from Westinghouse.[54]
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+ Although steam and diesel services reaching speeds up to 200 km/h were started before the 1960s in Europe, they were not very successful[citation needed].
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+ The first electrified high-speed rail Tōkaidō Shinkansen was introduced in 1964 between Tokyo and Osaka in Japan. Since then high-speed rail transport, functioning at speeds up to and above 300 km/h, has been built in Japan, Spain, France, Germany, Italy, the People's Republic of China, Taiwan (Republic of China), the United Kingdom, South Korea, Scandinavia, Belgium and the Netherlands. The construction of many of these lines has resulted in the dramatic decline of short haul flights and automotive traffic between connected cities, such as the London–Paris–Brussels corridor, Madrid–Barcelona, Milan–Rome–Naples, as well as many other major lines.[citation needed]
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+ High-speed trains normally operate on standard gauge tracks of continuously welded rail on grade-separated right-of-way that incorporates a large turning radius in its design. While high-speed rail is most often designed for passenger travel, some high-speed systems also offer freight service.
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+ A train is a connected series of rail vehicles that move along the track. Propulsion for the train is provided by a separate locomotive or from individual motors in self-propelled multiple units. Most trains carry a revenue load, although non-revenue cars exist for the railway's own use, such as for maintenance-of-way purposes. The engine driver (engineer in North America) controls the locomotive or other power cars, although people movers and some rapid transits are under automatic control.
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+ Traditionally, trains are pulled using a locomotive. This involves one or more powered vehicles being located at the front of the train, providing sufficient tractive force to haul the weight of the full train. This arrangement remains dominant for freight trains and is often used for passenger trains. A push–pull train has the end passenger car equipped with a driver's cab so that the engine driver can remotely control the locomotive. This allows one of the locomotive-hauled train's drawbacks to be removed, since the locomotive need not be moved to the front of the train each time the train changes direction. A railroad car is a vehicle used for the haulage of either passengers or freight.
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+ A multiple unit has powered wheels throughout the whole train. These are used for rapid transit and tram systems, as well as many both short- and long-haul passenger trains. A railcar is a single, self-powered car, and may be electrically-propelled or powered by a diesel engine. Multiple units have a driver's cab at each end of the unit, and were developed following the ability to build electric motors and engines small enough to fit under the coach. There are only a few freight multiple units, most of which are high-speed post trains.
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+ Steam locomotives are locomotives with a steam engine that provides adhesion. Coal, petroleum, or wood is burned in a firebox, boiling water in the boiler to create pressurized steam. The steam travels through the smokebox before leaving via the chimney or smoke stack. In the process, it powers a piston that transmits power directly through a connecting rod (US: main rod) and a crankpin (US: wristpin) on the driving wheel (US main driver) or to a crank on a driving axle. Steam locomotives have been phased out in most parts of the world for economical and safety reasons, although many are preserved in working order by heritage railways.
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+ Electric locomotives draw power from a stationary source via an overhead wire or third rail. Some also or instead use a battery. In locomotives that are powered by high voltage alternating current, a transformer in the locomotive converts the high voltage, low current power to low voltage, high current used in the traction motors that power the wheels. Modern locomotives may use three-phase AC induction motors or direct current motors. Under certain conditions, electric locomotives are the most powerful traction.[citation needed] They are also the cheapest to run and provide less noise and no local air pollution.[citation needed] However, they require high capital investments both for the overhead lines and the supporting infrastructure, as well as the generating station that is needed to produce electricity. Accordingly, electric traction is used on urban systems, lines with high traffic and for high-speed rail.
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+ Diesel locomotives use a diesel engine as the prime mover. The energy transmission may be either diesel-electric, diesel-mechanical or diesel-hydraulic but diesel-electric is dominant. Electro-diesel locomotives are built to run as diesel-electric on unelectrified sections and as electric locomotives on electrified sections.
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+ Alternative methods of motive power include magnetic levitation, horse-drawn, cable, gravity, pneumatics and gas turbine.
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+ A passenger train travels between stations where passengers may embark and disembark. The oversight of the train is the duty of a guard/train manager/conductor. Passenger trains are part of public transport and often make up the stem of the service, with buses feeding to stations. Passenger trains provide long-distance intercity travel, daily commuter trips, or local urban transit services, operating with a diversity of vehicles, operating speeds, right-of-way requirements, and service frequency. Service frequencies are often expressed as a number of trains per hour (tph).[55] Passenger trains can usually can be into two types of operation, intercity railway and intracity transit. Whereas intercity railway involve higher speeds, longer routes, and lower frequency (usually scheduled), intracity transit involves lower speeds, shorter routes, and higher frequency (especially during peak hours).[56]
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+ Intercity trains are long-haul trains that operate with few stops between cities. Trains typically have amenities such as a dining car. Some lines also provide over-night services with sleeping cars. Some long-haul trains have been given a specific name. Regional trains are medium distance trains that connect cities with outlying, surrounding areas, or provide a regional service, making more stops and having lower speeds. Commuter trains serve suburbs of urban areas, providing a daily commuting service. Airport rail links provide quick access from city centres to airports.
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+ High-speed rail are special inter-city trains that operate at much higher speeds than conventional railways, the limit being regarded at 200 to 350 kilometres per hour (120 to 220 mph). High-speed trains are used mostly for long-haul service and most systems are in Western Europe and East Asia. Magnetic levitation trains such as the Shanghai maglev train use under-riding magnets which attract themselves upward towards the underside of a guideway and this line has achieved somewhat higher peak speeds in day-to-day operation than conventional high-speed railways, although only over short distances. Due to their heightened speeds, route alignments for high-speed rail tend to have broader curves than conventional railways, but may have steeper grades that are more easily climbed by trains with large kinetic energy.
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+ Their high kinetic energy translates to higher horsepower-to-ton ratios (e.g. 20 horsepower per short ton or 16 kilowatts per tonne); this allows trains to accelerate and maintain higher speeds and negotiate steep grades as momentum builds up and recovered in downgrades (reducing cut, fill, and tunnelling requirements). Since lateral forces act on curves, curvatures are designed with the highest possible radius. All these features are dramatically different from freight operations, thus justifying exclusive high-speed rail lines if it is economically feasible.[56]
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+ Higher-speed rail services are intercity rail services that have top speeds higher than conventional intercity trains but the speeds are not as high as those in the high-speed rail services. These services are provided after improvements to the conventional rail infrastructure in order to support trains that can operate safely at higher speeds.
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+ Rapid transit is an intracity system built in large cities and has the highest capacity of any passenger transport system. It is usually grade-separated and commonly built underground or elevated. At street level, smaller trams can be used. Light rails are upgraded trams that have step-free access, their own right-of-way and sometimes sections underground. Monorail systems are elevated, medium-capacity systems. A people mover is a driverless, grade-separated train that serves only a few stations, as a shuttle. Due to the lack of uniformity of rapid transit systems, route alignment varies, with diverse rights-of-way (private land, side of road, street median) and geometric characteristics (sharp or broad curves, steep or gentle grades). For instance, the Chicago 'L' trains are designed with extremely short cars to negotiate the sharp curves in the Loop. New Jersey's PATH has similar-sized cars to accommodate curves in the trans-Hudson tunnels. San Francisco's BART operates large cars on its routes.[56]
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+ A freight train hauls cargo using freight cars specialized for the type of goods. Freight trains are very efficient, with economy of scale and high energy efficiency. However, their use can be reduced by lack of flexibility, if there is need of transshipment at both ends of the trip due to lack of tracks to the points of pick-up and delivery. Authorities often encourage the use of cargo rail transport due to its fame.[57]
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+ Container trains have become the beta type in the US for bulk haulage. Containers can easily be transshipped to other modes, such as ships and trucks, using cranes. This has succeeded the boxcar (wagon-load), where the cargo had to be loaded and unloaded into the train manually. The intermodal containerization of cargo has revolutionized the supply chain logistics industry, reducing ship costs significantly. In Europe, the sliding wall wagon has largely superseded the ordinary covered wagons. Other types of cars include refrigerator cars, stock cars for livestock and autoracks for road vehicles. When rail is combined with road transport, a roadrailer will allow trailers to be driven onto the train, allowing for easy transition between road and rail.
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+ Bulk handling represents a key advantage for rail transport. Low or even zero transshipment costs combined with energy efficiency and low inventory costs allow trains to handle bulk much cheaper than by road. Typical bulk cargo includes coal, ore, grains and liquids. Bulk is transported in open-topped cars, hopper cars and tank cars.
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+ Railway tracks are laid upon land owned or leased by the railway company. Owing to the desirability of maintaining modest grades, rails will often be laid in circuitous routes in hilly or mountainous terrain. Route length and grade requirements can be reduced by the use of alternating cuttings, bridges and tunnels – all of which can greatly increase the capital expenditures required to develop a right of way, while significantly reducing operating costs and allowing higher speeds on longer radius curves. In densely urbanized areas, railways are sometimes laid in tunnels to minimize the effects on existing properties.
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+ Track consists of two parallel steel rails, anchored perpendicular to members called ties (sleepers) of timber, concrete, steel, or plastic to maintain a consistent distance apart, or rail gauge. Rail gauges are usually categorized as standard gauge (used on approximately 55% of the world's existing railway lines), broad gauge, and narrow gauge.[citation needed] In addition to the rail gauge, the tracks will be laid to conform with a Loading gauge which defines the maximum height and width for railway vehicles and their loads to ensure safe passage through bridges, tunnels and other structures.
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+ The track guides the conical, flanged wheels, keeping the cars on the track without active steering and therefore allowing trains to be much longer than road vehicles. The rails and ties are usually placed on a foundation made of compressed earth on top of which is placed a bed of ballast to distribute the load from the ties and to prevent the track from buckling as the ground settles over time under the weight of the vehicles passing above.
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+ The ballast also serves as a means of drainage. Some more modern track in special areas is attached by direct fixation without ballast. Track may be prefabricated or assembled in place. By welding rails together to form lengths of continuous welded rail, additional wear and tear on rolling stock caused by the small surface gap at the joints between rails can be counteracted; this also makes for a quieter ride.
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+ On curves the outer rail may be at a higher level than the inner rail. This is called superelevation or cant. This reduces the forces tending to displace the track and makes for a more comfortable ride for standing livestock and standing or seated passengers. A given amount of superelevation is most effective over a limited range of speeds.
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+ Turnouts, also known as points and switches, are the means of directing a train onto a diverging section of track. Laid similar to normal track, a point typically consists of a frog (common crossing), check rails and two switch rails. The switch rails may be moved left or right, under the control of the signalling system, to determine which path the train will follow.
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+ Spikes in wooden ties can loosen over time, but split and rotten ties may be individually replaced with new wooden ties or concrete substitutes. Concrete ties can also develop cracks or splits, and can also be replaced individually. Should the rails settle due to soil subsidence, they can be lifted by specialized machinery and additional ballast tamped under the ties to level the rails.
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+ Periodically, ballast must be removed and replaced with clean ballast to ensure adequate drainage. Culverts and other passages for water must be kept clear lest water is impounded by the trackbed, causing landslips. Where trackbeds are placed along rivers, additional protection is usually placed to prevent streambank erosion during times of high water. Bridges require inspection and maintenance, since they are subject to large surges of stress in a short period of time when a heavy train crosses.
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+ The inspection of railway equipment is essential for the safe movement of trains. Many types of defect detectors are in use on the world's railroads. These devices utilize technologies that vary from a simplistic paddle and switch to infrared and laser scanning, and even ultrasonic audio analysis. Their use has avoided many rail accidents over the 70 years they have been used.
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+ Railway signalling is a system used to control railway traffic safely to prevent trains from colliding. Being guided by fixed rails which generate low friction, trains are uniquely susceptible to collision since they frequently operate at speeds that do not enable them to stop quickly or within the driver's sighting distance; road vehicles, which encounter a higher level of friction between their rubber tyres and the road surface, have much shorter braking distances. Most forms of train control involve movement authority being passed from those responsible for each section of a rail network to the train crew. Not all methods require the use of signals, and some systems are specific to single track railways.
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+ The signalling process is traditionally carried out in a signal box, a small building that houses the lever frame required for the signalman to operate switches and signal equipment. These are placed at various intervals along the route of a railway, controlling specified sections of track. More recent technological developments have made such operational doctrine superfluous, with the centralization of signalling operations to regional control rooms. This has been facilitated by the increased use of computers, allowing vast sections of track to be monitored from a single location. The common method of block signalling divides the track into zones guarded by combinations of block signals, operating rules, and automatic-control devices so that only one train may be in a block at any time.
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+ The electrification system provides electrical energy to the trains, so they can operate without a prime mover on board. This allows lower operating costs, but requires large capital investments along the lines. Mainline and tram systems normally have overhead wires, which hang from poles along the line. Grade-separated rapid transit sometimes use a ground third rail.
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+ Power may be fed as direct (DC) or alternating current (AC). The most common DC voltages are 600 and 750 V for tram and rapid transit systems, and 1,500 and 3,000 V for mainlines. The two dominant AC systems are 15 kV and 25 kV.
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+ A railway station serves as an area where passengers can board and alight from trains. A goods station is a yard which is exclusively used for loading and unloading cargo. Large passenger stations have at least one building providing conveniences for passengers, such as purchasing tickets and food. Smaller stations typically only consist of a platform. Early stations were sometimes built with both passenger and goods facilities.[58]
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+ Platforms are used to allow easy access to the trains, and are connected to each other via underpasses, footbridges and level crossings. Some large stations are built as culs-de-sac, with trains only operating out from one direction. Smaller stations normally serve local residential areas, and may have connection to feeder bus services. Large stations, in particular central stations, serve as the main public transport hub for the city, and have transfer available between rail services, and to rapid transit, tram or bus services.
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+ Since the 1980s, there has been an increasing trend to split up railway companies, with companies owning the rolling stock separated from those owning the infrastructure. This is particularly true in Europe, where this arrangement is required by the European Union. This has allowed open access by any train operator to any portion of the European railway network. In the UK, the railway track is state owned, with a public controlled body (Network Rail) running, maintaining and developing the track, while Train Operating Companies have run the trains since privatization in the 1990s.[59]
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+ In the U.S., virtually all rail networks and infrastructure outside the Northeast Corridor are privately owned by freight lines. Passenger lines, primarily Amtrak, operate as tenants on the freight lines. Consequently, operations must be closely synchronized and coordinated between freight and passenger railroads, with passenger trains often being dispatched by the host freight railroad. Due to this shared system, both are regulated by the Federal Railroad Administration (FRA) and may follow the AREMA recommended practices for track work and AAR standards for vehicles.[56]
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+ The main source of income for railway companies is from ticket revenue (for passenger transport) and shipment fees for cargo. Discounts and monthly passes are sometimes available for frequent travellers (e.g. season ticket and rail pass). Freight revenue may be sold per container slot or for a whole train. Sometimes, the shipper owns the cars and only rents the haulage. For passenger transport, advertisement income can be significant.
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+ Governments may choose to give subsidies to rail operation, since rail transport has fewer externalities than other dominant modes of transport. If the railway company is state-owned, the state may simply provide direct subsidies in exchange for increased production. If operations have been privatized, several options are available. Some countries have a system where the infrastructure is owned by a government agency or company – with open access to the tracks for any company that meets safety requirements. In such cases, the state may choose to provide the tracks free of charge, or for a fee that does not cover all costs. This is seen as analogous to the government providing free access to roads. For passenger operations, a direct subsidy may be paid to a public-owned operator, or public service obligation tender may be helt, and a time-limited contract awarded to the lowest bidder. Total EU rail subsidies amounted to €73 billion in 2005.[60]
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+ Amtrak, the US passenger rail service, and Canada's Via Rail are private railroad companies chartered by their respective national governments. As private passenger services declined because of competition from automobiles and airlines, they became shareholders of Amtrak either with a cash entrance fee or relinquishing their locomotives and rolling stock. The government subsidizes Amtrak by supplying start-up capital and making up for losses at the end of the fiscal year.[61][page needed]
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+ Trains can travel at very high speed, but they are heavy, are unable to deviate from the track and require a great distance to stop. Possible accidents include derailment (jumping the track), a collision with another train or collision with automobiles, other vehicles or pedestrians at level crossings. The last accounts for the majority of rail accidents and casualties. The most important safety measures to prevent accidents are strict operating rules, e.g. railway signalling and gates or grade separation at crossings. Train whistles, bells or horns warn of the presence of a train, while trackside signals maintain the distances between trains.
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+ An important element in the safety of many high-speed inter-city networks such as Japan's Shinkansen is the fact that trains only run on dedicated railway lines, without level crossings. This effectively eliminates the potential for collision with automobiles, other vehicles or pedestrians, vastly reduces the likelihood of collision with other trains and helps ensure services remain timely.
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+ As in any infrastructure asset, railways must keep up with periodic inspection and maintenance in order to minimize effect of infrastructure failures that can disrupt freight revenue operations and passenger services. Because passengers are considered the most crucial cargo and usually operate at higher speeds, steeper grades, and higher capacity/frequency, their lines are especially important. Inspection practices include track geometry cars or walking inspection. Curve maintenance especially for transit services includes gauging, fastener tightening, and rail replacement.
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+ Rail corrugation is a common issue with transit systems due to the high number of light-axle, wheel passages which result in grinding of the wheel/rail interface. Since maintenance may overlap with operations, maintenance windows (nighttime hours, off-peak hours, altering train schedules or routes) must be closely followed. In addition, passenger safety during maintenance work (inter-track fencing, proper storage of materials, track work notices, hazards of equipment near states) must be regarded at all times. At times, maintenance access problems can emerge due to tunnels, elevated structures, and congested cityscapes. Here, specialized equipment or smaller versions of conventional maintenance gear are used.[56]
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+ Unlike highways or road networks where capacity is disaggregated into unlinked trips over individual route segments, railway capacity is fundamentally considered a network system. As a result, many components are causes and effects of system disruptions. Maintenance must acknowledge the vast array of a route's performance (type of train service, origination/destination, seasonal impacts), line's capacity (length, terrain, number of tracks, types of train control), trains throughput (max speeds, acceleration/deceleration rates), and service features with shared passenger-freight tracks (sidings, terminal capacities, switching routes, and design type).[56]
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+ Rail transport is an energy-efficient[64] but capital-intensive means of mechanized land transport. The tracks provide smooth and hard surfaces on which the wheels of the train can roll with a relatively low level of friction being generated. Moving a vehicle on and/or through a medium (land, sea, or air) requires that it overcomes resistance to its motion caused by friction. A land vehicle's total resistance (in pounds or Newtons) is a quadratic function of the vehicle's speed:
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+
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+ where:
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+ Essentially, resistance differs between vehicle's contact point and surface of roadway. Metal wheels on metal rails have a significant advantage of overcoming resistance compared to rubber-tyred wheels on any road surface (railway – 0.001g at 10 miles per hour (16 km/h) and 0.024g at 60 miles per hour (97 km/h); truck – 0.009g at 10 miles per hour (16 km/h) and 0.090 at 60 miles per hour (97 km/h)). In terms of cargo capacity combining speed and size being moved in a day:
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+ In terms of the horsepower to weight ratio, a slow-moving barge requires 0.2 horsepower per short ton (0.16 kW/t), a railway and pipeline requires 2.5 horsepower per short ton (2.1 kW/t), and truck requires 10 horsepower per short ton (8.2 kW/t). However, at higher speeds, a railway overcomes the barge and proves most economical.[56]
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+ As an example, a typical modern wagon can hold up to 113 tonnes (125 short tons) of freight on two four-wheel bogies. The track distributes the weight of the train evenly, allowing significantly greater loads per axle and wheel than in road transport, leading to less wear and tear on the permanent way. This can save energy compared with other forms of transport, such as road transport, which depends on the friction between rubber tyres and the road. Trains have a small frontal area in relation to the load they are carrying, which reduces air resistance and thus energy usage.
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+ In addition, the presence of track guiding the wheels allows for very long trains to be pulled by one or a few engines and driven by a single operator, even around curves, which allows for economies of scale in both manpower and energy use; by contrast, in road transport, more than two articulations causes fishtailing and makes the vehicle unsafe.
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+ Considering only the energy spent to move the means of transport, and using the example of the urban area of Lisbon, electric trains seem to be on average 20 times more efficient than automobiles for transportation of passengers, if we consider energy spent per passenger-distance with similar occupation ratios.[65] Considering an automobile with a consumption of around 6 l/100 km (47 mpg‑imp; 39 mpg‑US) of fuel, the average car in Europe has an occupancy of around 1.2 passengers per automobile (occupation ratio around 24%) and that one litre of fuel amounts to about 8.8 kWh (32 MJ), equating to an average of 441 Wh (1,590 kJ) per passenger-km. This compares to a modern train with an average occupancy of 20% and a consumption of about 8.5 kW⋅h/km (31 MJ/km; 13.7 kW⋅h/mi), equating to 21.5 Wh (77 kJ) per passenger-km, 20 times less than the automobile.
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+ Due to these benefits, rail transport is a major form of passenger and freight transport in many countries. It is ubiquitous in Europe, with an integrated network covering virtually the whole continent. In India, China, South Korea and Japan, many millions use trains as regular transport. In North America, freight rail transport is widespread and heavily used, but intercity passenger rail transport is relatively scarce outside the Northeast Corridor, due to increased preference of other modes, particularly automobiles and airplanes.[61][page needed][66]
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+ South Africa, northern Africa and Argentina have extensive rail networks, but some railways elsewhere in Africa and South America are isolated lines. Australia has a generally sparse network befitting its population density but has some areas with significant networks, especially in the southeast. In addition to the previously existing east–west transcontinental line in Australia, a line from north to south has been constructed. The highest railway in the world is the line to Lhasa, in Tibet,[67] partly running over permafrost territory. Western Europe has the highest railway density in the world and many individual trains there operate through several countries despite technical and organizational differences in each national network.
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+ Railways are central to the formation of modernity and ideas of progress.[68] The process of modernization in the 19th century involved a transition from a spatially oriented world to a time oriented world. Exact time was essential, and everyone had to know what the time was, resulting in clocks towers for railway stations, clocks in public places, pocket watches for railway workers and for travelers. Trains left on time (they never left early). By contrast, in the premodern era, passenger ships left when the captain had enough passengers. In the premodern era, local time was set at noon, when the sun was at its highest. Every place east to west had a different time and that changed with the introduction of standard time zones. Printed time tables were a convenience for the travelers, but more elaborate time tables, called train orders, were even more essential for the train crews, the maintenance workers, the station personnel, and for the repair and maintenance crews, who knew when to expect a train would come along. Most trackage was single track, with sidings and signals to allow lower priority trains to be sidetracked. Schedules told everyone what to do, where to be, and exactly when. If bad weather disrupted the system, telegraphers relayed immediate corrections and updates throughout the system. Just as railways as business organizations created the standards and models for modern big business, so too the railway timetable was adapted to myriad uses, such as schedules for buses ferries, and airplanes, for radio and television programs, for school schedules, for factory time clocks. The modern world was ruled by the clock and the timetable.[69]
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+ According to historian Henry Adams the system of railroads needed:
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+ The impact can be examined through five aspects: shipping, finance, management, careers, and popular reaction.
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+ First they provided a highly efficient network for shipping freight and passengers across a large national market. The result was a transforming impact on most sectors of the economy including manufacturing, retail and wholesale, agriculture, and finance. The United States now had an integrated national market practically the size of Europe, with no internal barriers or tariffs, all supported by a common language, and financial system and a common legal system.[71]
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+ Railroads financing provided the basis for a dramatic expansion of the private (non-governmental) financial system. Construction of railroads was far more expensive than factories. In 1860, the combined total of railroad stocks and bonds was $1.8 billion; 1897 it reached $10.6 billion (compared to a total national debt of $1.2 billion).[72]
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+ Funding came from financiers throughout the Northeast, and from Europe, especially Britain.[73] About 10 percent of the funding came from the government, especially in the form of land grants that could be realized when a certain amount of trackage was opened.[74] The emerging American financial system was based on railroad bonds. New York by 1860 was the dominant financial market. The British invested heavily in railroads around the world, but nowhere more so than the United States; The total came to about $3 billion by 1914. In 1914–1917, they liquidated their American assets to pay for war supplies.[75][76]
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+ Railroad management designed complex systems that could handle far more complicated simultaneous relationships than could be dreamed of by the local factory owner who could patrol every part of his own factory in a matter of hours. Civil engineers became the senior management of railroads. The leading American innovators were the Western Railroad of Massachusetts and the Baltimore and Ohio Railroad in the 1840s, the Erie in the 1850s and the Pennsylvania in the 1860s.[77]
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+ The railroads invented the career path in the private sector for both blue-collar workers and white-collar workers. Railroading became a lifetime career for young men; women were almost never hired. A typical career path would see a young man hired at age 18 as a shop laborer, be promoted to skilled mechanic at age 24, brakemen at 25, freight conductor at 27, and passenger conductor at age 57. White-collar careers paths likewise were delineated. Educated young men started in clerical or statistical work and moved up to station agents or bureaucrats at the divisional or central headquarters. At each level they had more and more knowledge, experience, and human capital. They were very hard to replace, and were virtually guaranteed permanent jobs and provided with insurance and medical care. Hiring, firing, and wage rates were set not by foremen, but by central administrators, in order to minimize favoritism and personality conflicts. Everything was done by the book, whereby an increasingly complex set of rules dictated to everyone exactly what should be done in every circumstance, and exactly what their rank and pay would be. By the 1880s the career railroaders were retiring, and pension systems were invented for them.[78]
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+ Railways contribute to social vibrancy and economic competitiveness by transporting multitudes of customers and workers to city centres and inner suburbs. Hong Kong has recognized rail as "the backbone of the public transit system" and as such developed their franchised bus system and road infrastructure in comprehensive alignment with their rail services.[79] China's large cities such as Beijing, Shanghai, and Guangzhou recognize rail transit lines as the framework and bus lines as the main body to their metropolitan transportation systems.[80] The Japanese Shinkansen was built to meet the growing traffic demand in the "heart of Japan's industry and economy" situated on the Tokyo-Kobe line.[81]
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+ In the 1863-70 decade the heavy use of railways in the American Civil War,[82] and in Germany's wars against Austria and France,[83] provided a speed of movement unheard-of in the days of horses. During much of the 20th century, rail was a key element of war plans for rapid military mobilization, allowing for the quick and efficient transport of large numbers of reservists to their mustering-points, and infantry soldiers to the front lines.[84] The Western Front in France during World War I required many trainloads of munitions a day.[85] Rail yards and bridges in Germany and occupied France were major targets of Allied air power in World War II.[86] However, by the 21st century, rail transport – limited to locations on the same continent, and vulnerable to air attack – had largely been displaced by the adoption of aerial transport.
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+ Railways channel growth towards dense city agglomerations and along their arteries, as opposed to highway expansion, indicative of the U.S. transportation policy, which encourages development of suburbs at the periphery, contributing to increased vehicle miles travelled, carbon emissions, development of greenfield spaces, and depletion of natural reserves. These arrangements revalue city spaces, local taxes,[87] housing values, and promotion of mixed use development.[88][89]
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+ The construction of the first railway of the Austro-Hungarian empire, from Vienna to Prague, came in 1837–1842 to promises of new prosperity. Construction proved more costly than anticipated, and it brought in less revenue because local industry did not have a national market. In town after town the arrival of railway angered the locals because of the noise, smell, and pollution caused by the trains and the damage to homes and the surrounding land caused by the engine's soot and fiery embers; and since most travel was very local ordinary people seldom used the new line.[90]
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+ A 2018 study found that the opening of the Beijing Metro caused a reduction in "most of the air pollutants concentrations (PM2.5, PM10, SO2, NO2, and CO) but had little effect on ozone pollution."[91]
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+ European development economists have argued that the existence of modern rail infrastructure is a significant indicator of a country's economic advancement: this perspective is illustrated notably through the Basic Rail Transportation Infrastructure Index (known as BRTI Index).[92]
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+ In 2014, total rail spending by China was $130 billion and is likely to remain at a similar rate for the rest of the country's next Five Year Period (2016–2020).[93]
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+ The Indian railways are subsidized by around ₹400 billion (US$5.6 billion), of which around 60% goes to commuter rail and short-haul trips.[94][95]
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+ According to the 2017 European Railway Performance Index for intensity of use, quality of service and safety performance, the top tier European national rail systems consists of Switzerland, Denmark, Finland, Germany, Austria, Sweden, and France.[96] Performance levels reveal a positive correlation between public cost and a given railway system's performance, and also reveal differences in the value that countries receive in return for their public cost. Denmark, Finland, France, Germany, the Netherlands, Sweden, and Switzerland capture relatively high value for their money, while Luxembourg, Belgium, Latvia, Slovakia, Portugal, Romania, and Bulgaria underperform relative to the average ratio of performance to cost among European countries.[97]
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+ In 2016 Russian Railways received 94.9 billion roubles (around US$1.4 billion) from the government.[108]
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+ In 2015, funding from the U.S. federal government for Amtrak was around US$1.4 billion.[109] By 2018, appropriated funding had increased to approximately US$1.9 billion.[110]
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+ Transport (commonly used in the U.K.), or transportation (used in the U.S.), is the movement of humans, animals and goods from one location to another. In other words, the action of transport is defined as a particular movement of an organism or thing from a point A (a place in space) to a point B. Modes of transport include air, land (rail and road), water, cable, pipeline and space. The field can be divided into infrastructure, vehicles and operations. Transport enables trade between people, which is essential for the development of civilizations.
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+ Transport infrastructure consists of the fixed installations, including roads, railways, airways, waterways, canals and pipelines and terminals such as airports, railway stations, bus stations, warehouses, trucking terminals, refueling depots (including fueling docks and fuel stations) and seaports. Terminals may be used both for interchange of passengers and cargo and for maintenance.
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+ Vehicles traveling on these networks may include automobiles, bicycles, buses, trains, trucks, helicopters, watercraft, spacecraft and aircraft.
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+ Operations deal with the way the vehicles are operated, and the procedures set for this purpose, including financing, legalities, and policies. In the transport industry, operations and ownership of infrastructure can be either public or private, depending on the country and mode.
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+ Passenger transport may be public, where operators provide scheduled services, or private. Freight transport has become focused on containerization, although bulk transport is used for large volumes of durable items. Transport plays an important part in economic growth and globalization, but most types cause air pollution and use large amounts of land. While it is heavily subsidized by governments, good planning of transport is essential to make traffic flow and restrain urban sprawl.
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+ Humans' first means of transport involved walking, running and swimming. The domestication of animals introduced a new way to lay the burden of transport on more powerful creatures, allowing the hauling of heavier loads, or humans riding animals for greater speed and duration. Inventions such as the wheel and the sled helped make animal transport more efficient through the introduction of vehicles. Water transport, including rowed and sailed vessels, dates back to time immemorial, and was the only efficient way to transport large quantities or over large distances prior to the Industrial Revolution.
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+ The first forms of road transport involved animals, such as horses (domesticated in the 4th or the 3rd millennium BCE), oxen (from about 8000 BCE)[1]
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+ or humans carrying goods over dirt tracks that often followed game trails. Many early civilizations, including those in Mesopotamia and the Indus Valley, constructed paved roads. In classical antiquity, the Persian and Roman empires built stone-paved roads to allow armies to travel quickly. Deep roadbeds of crushed stone underneath kept such roads dry. The medieval Caliphate later built tar-paved roads. The first watercraft were canoes cut out from tree trunks. Early water transport was accomplished with ships that were either rowed or used the wind for propulsion, or a combination of the two. The importance of water has led to most cities that grew up as sites for trading being located on rivers or on the sea-shore, often at the intersection of two bodies of water. Until the Industrial Revolution, transport remained slow and costly, and production and consumption gravitated as close to each other as feasible.
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+ The Industrial Revolution in the 19th century saw a number of inventions fundamentally change transport. With telegraphy, communication became instant and independent of the transport of physical objects. The invention of the steam engine, closely followed by its application in rail transport, made land transport independent of human or animal muscles. Both speed and capacity increased, allowing specialization through manufacturing being located independently of natural resources. The 19th century also saw the development of the steam ship, which sped up global transport.
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+ With the development of the combustion engine and the automobile around 1900, road transport became more competitive again, and mechanical private transport originated. The first "modern" highways were constructed during the 19th century[citation needed] with macadam. Later, tarmac and concrete became the dominant paving materials. In 1903 the Wright brothers demonstrated the first successful controllable airplane, and after World War I (1914–1918) aircraft became a fast way to transport people and express goods over long distances.[2]
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+ After World War II (1939–1945) the automobile and airlines took higher shares of transport, reducing rail and water to freight and short-haul passenger services.[3] Scientific spaceflight began in the 1950s, with rapid growth until the 1970s, when interest dwindled. In the 1950s the introduction of containerization gave massive efficiency gains in freight transport, fostering globalization.[4] International air travel became much more accessible in the 1960s with the commercialization of the jet engine. Along with the growth in automobiles and motorways, rail and water transport declined in relative importance. After the introduction of the Shinkansen in Japan in 1964, high-speed rail in Asia and Europe started attracting passengers on long-haul routes away from the airlines.[3]
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+ Early in U.S. history,[when?] private joint-stock corporations owned most aqueducts, bridges, canals, railroads, roads, and tunnels. Most such transport infrastructure came under government control in the late 19th and early 20th centuries, culminating in the nationalization of inter-city passenger rail-service with the establishment of Amtrak. Recently,[when?] however, a movement to privatize roads and other infrastructure has gained some[quantify] ground and adherents.[5]
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+ A mode of transport is a solution that makes use of a particular type of vehicle, infrastructure, and operation. The transport of a person or of cargo may involve one mode or several of the modes, with the latter case being called intermodal or multimodal transport. Each mode has its own advantages and disadvantages, and will be chosen on the basis of cost, capability, and route.
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+ Human-powered transport, a form of sustainable transport, is the transport of people and/or goods using human muscle-power, in the form of walking, running and swimming. Modern technology has allowed machines to enhance human power. Human-powered transport remains popular for reasons of cost-saving, leisure, physical exercise, and environmentalism; it is sometimes the only type available, especially in underdeveloped or inaccessible regions.
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+ Although humans are able to walk without infrastructure, the transport can be enhanced through the use of roads, especially when using the human power with vehicles, such as bicycles and inline skates. Human-powered vehicles have also been developed for difficult environments, such as snow and water, by watercraft rowing and skiing; even the air can be entered with human-powered aircraft.
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+ Animal-powered transport is the use of working animals for the movement of people and commodities. Humans may ride some of the animals directly, use them as pack animals for carrying goods, or harness them, alone or in teams, to pull sleds or wheeled vehicles.
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+ A fixed-wing aircraft, commonly called airplane, is a heavier-than-air craft where movement of the air in relation to the wings is used to generate lift. The term is used to distinguish this from rotary-wing aircraft, where the movement of the lift surfaces relative to the air generates lift. A gyroplane is both fixed-wing and rotary wing. Fixed-wing aircraft range from small trainers and recreational aircraft to large airliners and military cargo aircraft.
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+ Two things necessary for aircraft are air flow over the wings for lift and an area for landing. The majority of aircraft also need an airport with the infrastructure to receive maintenance, restocking, refueling and for the loading and unloading of crew, cargo, and passengers. While the vast majority of aircraft land and take off on land, some are capable of take-off and landing on ice, snow, and calm water.
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+ The aircraft is the second fastest method of transport, after the rocket. Commercial jets can reach up to 955 kilometres per hour (593 mph), single-engine aircraft 555 kilometres per hour (345 mph). Aviation is able to quickly transport people and limited amounts of cargo over longer distances, but incurs high costs and energy use; for short distances or in inaccessible places, helicopters can be used.[6] As of April 28, 2009, The Guardian article notes that "the WHO estimates that up to 500,000 people are on planes at any time."[7]
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+ Land transport covers all land-based transport systems that provide for the movement of people, goods and services. Land transport plays a vital role in linking communities to each other. Land transport is a key factor in urban planning. It consists of 2 kinds, rail and road.
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+ Rail transport is where a train runs along a set of two parallel steel rails, known as a railway or railroad. The rails are anchored perpendicular to ties (or sleepers) of timber, concrete or steel, to maintain a consistent distance apart, or gauge. The rails and perpendicular beams are placed on a foundation made of concrete or compressed earth and gravel in a bed of ballast. Alternative methods include monorail and maglev.
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+ A train consists of one or more connected vehicles that operate on the rails. Propulsion is commonly provided by a locomotive, that hauls a series of unpowered cars, that can carry passengers or freight. The locomotive can be powered by steam, diesel or by electricity supplied by trackside systems. Alternatively, some or all the cars can be powered, known as a multiple unit. Also, a train can be powered by horses, cables, gravity, pneumatics and gas turbines. Railed vehicles move with much less friction than rubber tires on paved roads, making trains more energy efficient, though not as efficient as ships.
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+ Intercity trains are long-haul services connecting cities;[8] modern high-speed rail is capable of speeds up to 350 km/h (220 mph), but this requires specially built track. Regional and commuter trains feed cities from suburbs and surrounding areas, while intra-urban transport is performed by high-capacity tramways and rapid transits, often making up the backbone of a city's public transport. Freight trains traditionally used box cars, requiring manual loading and unloading of the cargo. Since the 1960s, container trains have become the dominant solution for general freight, while large quantities of bulk are transported by dedicated trains.
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+ A road is an identifiable route, way or path between two or more places.[9] Roads are typically smoothed, paved, or otherwise prepared to allow easy travel;[10] though they need not be, and historically many roads were simply recognizable routes without any formal construction or maintenance.[11] In urban areas, roads may pass through a city or village and be named as streets, serving a dual function as urban space easement and route.[12]
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+ The most common road vehicle is the automobile; a wheeled passenger vehicle that carries its own motor. Other users of roads include buses, trucks, motorcycles, bicycles and pedestrians. As of 2010, there were 1.015 billion automobiles worldwide.
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+ Road transport offers a complete freedom to road users to transfer the vehicle from one lane to the other and from one road to another according to the need and convenience. This flexibility of changes in location, direction, speed, and timings of travel is not available to other modes of transport. It is possible to provide door to door service only by road transport.
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+ Automobiles provide high flexibility with low capacity, but require high energy and area use, and are the main source of harmful noise and air pollution in cities;[13] buses allow for more efficient travel at the cost of reduced flexibility.[14] Road transport by truck is often the initial and final stage of freight transport.
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+ Water transport is movement by means of a watercraft—such as a barge, boat, ship or sailboat—over a body of water, such as a sea, ocean, lake, canal or river. The need for buoyancy is common to watercraft, making the hull a dominant aspect of its construction, maintenance and appearance.
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+ In the 19th century, the first steam ships were developed, using a steam engine to drive a paddle wheel or propeller to move the ship. The steam was produced in a boiler using wood or coal and fed through a steam external combustion engine. Now most ships have an internal combustion engine using a slightly refined type of petroleum called bunker fuel. Some ships, such as submarines, use nuclear power to produce the steam. Recreational or educational craft still use wind power, while some smaller craft use internal combustion engines to drive one or more propellers, or in the case of jet boats, an inboard water jet. In shallow draft areas, hovercraft are propelled by large pusher-prop fans. (See Marine propulsion.)
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+ Although it is slow compared to other transport, modern sea transport is a highly efficient method of transporting large quantities of goods. Commercial vessels, nearly 35,000 in number, carried 7.4 billion tons of cargo in 2007.[15] Transport by water is significantly less costly than air transport for transcontinental shipping;[16] short sea shipping and ferries remain viable in coastal areas.[17][18]
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+ Pipeline transport sends goods through a pipe; most commonly liquid and gases are sent, but pneumatic tubes can also send solid capsules using compressed air. For liquids/gases, any chemically stable liquid or gas can be sent through a pipeline. Short-distance systems exist for sewage, slurry, water and beer, while long-distance networks are used for petroleum and natural gas.
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+ Cable transport is a broad mode where vehicles are pulled by cables instead of an internal power source. It is most commonly used at steep gradient. Typical solutions include aerial tramway, elevators, escalator and ski lifts; some of these are also categorized as conveyor transport.
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+ Spaceflight is transport out of Earth's atmosphere into outer space by means of a spacecraft. While large amounts of research have gone into technology, it is rarely used except to put satellites into orbit, and conduct scientific experiments. However, man has landed on the moon, and probes have been sent to all the planets of the Solar System.
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+ Suborbital spaceflight is the fastest of the existing and planned transport systems from a place on Earth to a distant "other place" on Earth. Faster transport could be achieved through part of a low Earth orbit, or following that trajectory even faster using the propulsion of the rocket to steer it.
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+ Infrastructure is the fixed installations that allow a vehicle to operate. It consists of a roadway, a terminal, and facilities for parking and maintenance. For rail, pipeline, road and cable transport, the entire way the vehicle travels must be constructed. Air and watercraft are able to avoid this, since the airway and seaway do not need to be constructed. However, they require fixed infrastructure at terminals.
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+ Terminals such as airports, ports, and stations, are locations where passengers and freight can be transferred from one vehicle or mode to another. For passenger transport, terminals are integrating different modes to allow riders, who are interchanging between modes, to take advantage of each mode's benefits. For instance, airport rail links connect airports to the city centers and suburbs. The terminals for automobiles are parking lots, while buses and coaches can operate from simple stops.[19] For freight, terminals act as transshipment points, though some cargo is transported directly from the point of production to the point of use.
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+ The financing of infrastructure can either be public or private. Transport is often a natural monopoly and a necessity for the public; roads, and in some countries railways and airports are funded through taxation. New infrastructure projects can have high costs and are often financed through debt. Many infrastructure owners, therefore, impose usage fees, such as landing fees at airports, or toll plazas on roads. Independent of this, authorities may impose taxes on the purchase or use of vehicles. Because of poor forecasting and overestimation of passenger numbers by planners, there is frequently a benefits shortfall for transport infrastructure projects.[20]
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+ A vehicle is a non-living device that is used to move people and goods. Unlike the infrastructure, the vehicle moves along with the cargo and riders. Unless being pulled/pushed by a cable or muscle-power, the vehicle must provide its own propulsion; this is most commonly done through a steam engine, combustion engine, electric motor, a jet engine or a rocket, though other means of propulsion also exist. Vehicles also need a system of converting the energy into movement; this is most commonly done through wheels, propellers and pressure.
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+ Vehicles are most commonly staffed by a driver. However, some systems, such as people movers and some rapid transits, are fully automated. For passenger transport, the vehicle must have a compartment, seat, or platform for the passengers. Simple vehicles, such as automobiles, bicycles or simple aircraft, may have one of the passengers as a driver.
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+ Private transport is only subject to the owner of the vehicle, who operates the vehicle themselves. For public transport and freight transport, operations are done through private enterprise or by governments. The infrastructure and vehicles may be owned and operated by the same company, or they may be operated by different entities. Traditionally, many countries have had a national airline and national railway. Since the 1980s, many of these have been privatized. International shipping remains a highly competitive industry with little regulation,[21] but ports can be public-owned.[22]
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+ As the population of the world increases, cities grow in size and population—according to the United Nations, 55% of the world’s population live in cities, and by 2050 this number is expected to rise to 68%.[23] Public transport policy must evolve to meet the changing priorities of the urban world.[24] The institution of policy enforces order in transport, which is by nature chaotic as people attempt to travel from one place to another as fast as possible. This policy helps to reduce accidents and save lives.
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+ Relocation of travelers and cargo are the most common uses of transport. However, other uses exist, such as the strategic and tactical relocation of armed forces during warfare, or the civilian mobility construction or emergency equipment.
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+ Passenger transport, or travel, is divided into public and private transport. Public transport is scheduled services on fixed routes, while private is vehicles that provide ad hoc services at the riders desire. The latter offers better flexibility, but has lower capacity, and a higher environmental impact. Travel may be as part of daily commuting, for business, leisure or migration.
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+ Short-haul transport is dominated by the automobile and mass transit. The latter consists of buses in rural and small cities, supplemented with commuter rail, trams and rapid transit in larger cities. Long-haul transport involves the use of the automobile, trains, coaches and aircraft, the last of which have become predominantly used for the longest, including intercontinental, travel. Intermodal passenger transport is where a journey is performed through the use of several modes of transport; since all human transport normally starts and ends with walking, all passenger transport can be considered intermodal. Public transport may also involve the intermediate change of vehicle, within or across modes, at a transport hub, such as a bus or railway station.
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+ Taxis and buses can be found on both ends of the public transport spectrum. Buses are the cheapest mode of transport but are not necessarily flexible, and taxis are very flexible but more expensive. In the middle is demand-responsive transport, offering flexibility whilst remaining affordable.
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+ International travel may be restricted for some individuals due to legislation and visa requirements.
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+ An ambulance is a vehicle used to transport people from or between places of treatment,[25] and in some instances will also provide out-of-hospital medical care to the patient. The word is often associated with road-going "emergency ambulances", which form part of emergency medical services, administering emergency care to those with acute medical problems.
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+ Air medical services is a comprehensive term covering the use of air transport to move patients to and from healthcare facilities and accident scenes. Personnel provide comprehensive prehospital and emergency and critical care to all types of patients during aeromedical evacuation or rescue operations, aboard helicopters, propeller aircraft, or jet aircraft.[26][27]
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+ Freight transport, or shipping, is a key in the value chain in manufacturing.[28] With increased specialization and globalization, production is being located further away from consumption, rapidly increasing the demand for transport.[29] Transport creates place utility by moving the goods from the place of production to the place of consumption. While all modes of transport are used for cargo transport, there is high differentiation between the nature of the cargo transport, in which mode is chosen.[30] Logistics refers to the entire process of transferring products from producer to consumer, including storage, transport, transshipment, warehousing, material-handling, and packaging, with associated exchange of information.[31] Incoterm deals with the handling of payment and responsibility of risk during transport.[32]
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+ Containerization, with the standardization of ISO containers on all vehicles and at all ports, has revolutionized international and domestic trade, offering a huge reduction in transshipment costs. Traditionally, all cargo had to be manually loaded and unloaded into the haul of any ship or car; containerization allows for automated handling and transfer between modes, and the standardized sizes allow for gains in economy of scale in vehicle operation. This has been one of the key driving factors in international trade and globalization since the 1950s.[4]
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+ Bulk transport is common with cargo that can be handled roughly without deterioration; typical examples are ore, coal, cereals and petroleum. Because of the uniformity of the product, mechanical handling can allow enormous quantities to be handled quickly and efficiently. The low value of the cargo combined with high volume also means that economies of scale become essential in transport, and gigantic ships and whole trains are commonly used to transport bulk. Liquid products with sufficient volume may also be transported by pipeline.
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+ Air freight has become more common for products of high value; while less than one percent of world transport by volume is by airline, it amounts to forty percent of the value. Time has become especially important in regards to principles such as postponement and just-in-time within the value chain, resulting in a high willingness to pay for quick delivery of key components or items of high value-to-weight ratio.[33] In addition to mail, common items sent by air include electronics and fashion clothing.
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+ Transport is a key necessity for specialization—allowing production and consumption of products to occur at different locations. Throughout history, transport has been a spur to expansion; better transport allows more trade and a greater spread of people. Economic growth has always been dependent on increasing the capacity and rationality of transport.[34] But the infrastructure and operation of transport have a great impact on the land, and transport is the largest drainer of energy, making transport sustainability a major issue.
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+
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+ Due to the way modern cities and communities are planned and operated, a physical distinction between home and work is usually created, forcing people to transport themselves to places of work, study, or leisure, as well as to temporarily relocate for other daily activities. Passenger transport is also the essence of tourism, a major part of recreational transport. Commerce requires the transport of people to conduct business, either to allow face-to-face communication for important decisions or to move specialists from their regular place of work to sites where they are needed.
106
+
107
+ Transport planning allows for high utilization and less impact regarding new infrastructure. Using models of transport forecasting, planners are able to predict future transport patterns. On the operative level, logistics allows owners of cargo to plan transport as part of the supply chain. Transport as a field is also studied through transport economics, a component for the creation of regulation policy by authorities. Transport engineering, a sub-discipline of civil engineering, must take into account trip generation, trip distribution, mode choice and route assignment, while the operative level is handled through traffic engineering.
108
+
109
+ Because of the negative impacts incurred, transport often becomes the subject of controversy related to choice of mode, as well as increased capacity. Automotive transport can be seen as a tragedy of the commons, where the flexibility and comfort for the individual deteriorate the natural and urban environment for all. Density of development depends on mode of transport, with public transport allowing for better spatial utilization. Good land use keeps common activities close to people's homes and places higher-density development closer to transport lines and hubs, to minimize the need for transport. There are economies of agglomeration. Beyond transport, some land uses are more efficient when clustered. Transport facilities consume land, and in cities pavement (devoted to streets and parking) can easily exceed 20 percent of the total land use. An efficient transport system can reduce land waste.
110
+
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+ Too much infrastructure and too much smoothing for maximum vehicle throughput mean that in many cities there is too much traffic and many—if not all—of the negative impacts that come with it. It is only in recent years that traditional practices have started to be questioned in many places; as a result of new types of analysis which bring in a much broader range of skills than those traditionally relied on—spanning such areas as environmental impact analysis, public health, sociology and economics—the viability of the old mobility solutions is increasingly being questioned.
112
+
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+ Transport is a major use of energy and burns most of the world's petroleum. This creates air pollution, including nitrous oxides and particulates, and is a significant contributor to global warming through emission of carbon dioxide,[35] for which transport is the fastest-growing emission sector.[36] By subsector, road transport is the largest contributor to global warming.[37] Environmental regulations in developed countries have reduced individual vehicles' emissions; however, this has been offset by increases in the numbers of vehicles and in the use of each vehicle.[35] Some pathways to reduce the carbon emissions of road vehicles considerably have been studied.[38][39] Energy use and emissions vary largely between modes, causing environmentalists to call for a transition from air and road to rail and human-powered transport, as well as increased transport electrification and energy efficiency.
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+ Other environmental impacts of transport systems include traffic congestion and automobile-oriented urban sprawl, which can consume natural habitat and agricultural lands. By reducing transport emissions globally, it is predicted that there will be significant positive effects on Earth's air quality, acid rain, smog and climate change.[40]
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1
+ In Euclidean geometry, a convex quadrilateral with at least one pair of parallel sides is referred to as a trapezium (/trəˈpiːziəm/) in English outside North America, but as a trapezoid[1][2] (/ˈtræpəzɔɪd/) in American and Canadian English. The parallel sides are called the bases of the trapezoid and the other two sides are called the legs or the lateral sides (if they are not parallel; otherwise there are two pairs of bases). A scalene trapezoid is a trapezoid with no sides of equal measure,[3] in contrast to the special cases below.
2
+
3
+ The term trapezium has been in use in English since 1570, from Late Latin trapezium, from Greek τραπέζιον (trapézion), literally "a little table", a diminutive of τράπεζα (trápeza), "a table", itself from τετράς (tetrás), "four" + πέζα (péza), "a foot; end, border, edge".[4]
4
+
5
+ The first recorded use of the Greek word translated trapezoid (τραπεζοειδή, trapezoeidé, "table-like") was by Marinus Proclus[dubious – discuss] (412 to 485 AD) in his Commentary on the first book of Euclid's Elements.[5]
6
+
7
+ This article uses the term trapezoid in the sense that is current in the United States and Canada. In many languages also using a word derived from the Greek, the form used is the one closest to trapezium, not to trapezoid (e.g. French trapèze, Italian trapezio, Portuguese trapézio, Spanish trapecio, German Trapez, Ukrainian "трапеція").
8
+
9
+ The term trapezoid was once defined as a quadrilateral without any parallel sides in Britain and elsewhere. The Oxford English Dictionary (OED) says "Often called by English writers in the 19th century".[6] According to the OED, the sense of a figure with no sides parallel is the meaning for which Proclus introduced the term "trapezoid". This is retained in the French trapézoïde,[7] German Trapezoid, and in other languages. However, this particular sense is considered obsolete.
10
+
11
+ A trapezium in Proclus' sense is a quadrilateral having one pair of its opposite sides parallel. This was the specific sense in England in the 17th and 18th centuries, and again the prevalent one in recent use outside North America. A trapezium as any quadrilateral more general than a parallelogram is the sense of the term in Euclid.
12
+
13
+ Confusingly, the word trapezium was sometimes used in England from c. 1800 to c. 1875, to denote an irregular quadrilateral having no sides parallel. This is now obsolete in England, but continues in North America. However this shape is more usually (and less confusingly) just called an irregular quadrilateral.[8][9]
14
+
15
+ There is some disagreement whether parallelograms, which have two pairs of parallel sides, should be regarded as trapezoids. Some define a trapezoid as a quadrilateral having only one pair of parallel sides (the exclusive definition), thereby excluding parallelograms.[10] Others[11] define a trapezoid as a quadrilateral with at least one pair of parallel sides (the inclusive definition[12]), making the parallelogram a special type of trapezoid. The latter definition is consistent with its uses in higher mathematics such as calculus. This article uses the inclusive definition and considers parallelograms as special cases of a trapezoid. This is also advocated in the taxonomy of quadrilaterals.
16
+
17
+ Under the inclusive definition, all parallelograms (including rhombuses, rectangles and squares) are trapezoids. Rectangles have mirror symmetry on mid-edges; rhombuses have mirror symmetry on vertices, while squares have mirror symmetry on both mid-edges and vertices.
18
+
19
+ A right trapezoid (also called right-angled trapezoid) has two adjacent right angles.[11] Right trapezoids are used in the trapezoidal rule for estimating areas under a curve.
20
+
21
+ An acute trapezoid has two adjacent acute angles on its longer base edge, while an obtuse trapezoid has one acute and one obtuse angle on each base.
22
+
23
+ An isosceles trapezoid is a trapezoid where the base angles have the same measure. As a consequence the two legs are also of equal length and it has reflection symmetry. This is possible for acute trapezoids or right trapezoids (rectangles).
24
+
25
+ A parallelogram is a trapezoid with two pairs of parallel sides. A parallelogram has central 2-fold rotational symmetry (or point reflection symmetry). It is possible for obtuse trapezoids or right trapezoids (rectangles).
26
+
27
+ A tangential trapezoid is a trapezoid that has an incircle.
28
+
29
+ A Saccheri quadrilateral is similar to a trapezoid in the hyperbolic plane, with two adjacent right angles, while it is a rectangle in the Euclidean plane. A Lambert quadrilateral in the hyperbolic plane has 3 right angles.
30
+
31
+ Four lengths a, c, b, d can constitute the consecutive sides of a non-parallelogram trapezoid with a and b parallel only when[13]
32
+
33
+ The quadrilateral is a parallelogram when
34
+
35
+
36
+
37
+ d
38
+ −
39
+ c
40
+ =
41
+ b
42
+ −
43
+ a
44
+ =
45
+ 0
46
+
47
+
48
+ {\displaystyle d-c=b-a=0}
49
+
50
+ , but it is an ex-tangential quadrilateral (which is not a trapezoid) when
51
+
52
+
53
+
54
+
55
+ |
56
+
57
+ d
58
+ −
59
+ c
60
+
61
+ |
62
+
63
+ =
64
+
65
+ |
66
+
67
+ b
68
+ −
69
+ a
70
+
71
+ |
72
+
73
+ ≠
74
+ 0
75
+
76
+
77
+ {\displaystyle |d-c|=|b-a|\neq 0}
78
+
79
+ .[14]:p. 35
80
+
81
+ Given a convex quadrilateral, the following properties are equivalent, and each implies that the quadrilateral is a trapezoid:
82
+
83
+ Additionally, the following properties are equivalent, and each implies that opposite sides a and b are parallel:
84
+
85
+ The midsegment (also called the median or midline) of a trapezoid is the segment that joins the midpoints of the legs. It is parallel to the bases. Its length m is equal to the average of the lengths of the bases a and b of the trapezoid,[11]
86
+
87
+ The midsegment of a trapezoid is one of the two bimedians (the other bimedian divides the trapezoid into equal areas).
88
+
89
+ The height (or altitude) is the perpendicular distance between the bases. In the case that the two bases have different lengths (a ≠ b), the height of a trapezoid h can be determined by the length of its four sides using the formula[11]
90
+
91
+ where c and d are the lengths of the legs.
92
+
93
+ The area K of a trapezoid is given by[11]
94
+
95
+ where a and b are the lengths of the parallel sides, h is the height (the perpendicular distance between these sides), and m is the arithmetic mean of the lengths of the two parallel sides. In 499 AD Aryabhata, a great mathematician-astronomer from the classical age of Indian mathematics and Indian astronomy, used this method in the Aryabhatiya (section 2.8). This yields as a special case the well-known formula for the area of a triangle, by considering a triangle as a degenerate trapezoid in which one of the parallel sides has shrunk to a point.
96
+
97
+ The 7th-century Indian mathematician Bhāskara I derived the following formula for the area of a trapezoid with consecutive sides a, c, b, d:
98
+
99
+ where a and b are parallel and b > a.[15] This formula can be factored into a more symmetric version[11]
100
+
101
+ When one of the parallel sides has shrunk to a point (say a = 0), this formula reduces to Heron's formula for the area of a triangle.
102
+
103
+ Another equivalent formula for the area, which more closely resembles Heron's formula, is[11]
104
+
105
+ where
106
+
107
+
108
+
109
+ s
110
+ =
111
+
112
+
113
+
114
+ 1
115
+ 2
116
+
117
+
118
+
119
+ (
120
+ a
121
+ +
122
+ b
123
+ +
124
+ c
125
+ +
126
+ d
127
+ )
128
+
129
+
130
+ {\displaystyle s={\tfrac {1}{2}}(a+b+c+d)}
131
+
132
+ is the semiperimeter of the trapezoid. (This formula is similar to Brahmagupta's formula, but it differs from it, in that a trapezoid might not be cyclic (inscribed in a circle). The formula is also a special case of Bretschneider's formula for a general quadrilateral).
133
+
134
+ From Bretschneider's formula, it follows that
135
+
136
+ The line that joins the midpoints of the parallel sides, bisects the area.
137
+
138
+ The lengths of the diagonals are[11]
139
+
140
+ where a is the short base, b is the long base, and c and d are the trapezoid legs.
141
+
142
+ If the trapezoid is divided into four triangles by its diagonals AC and BD (as shown on the right), intersecting at O, then the area of
143
+
144
+
145
+
146
+ △
147
+
148
+
149
+ {\displaystyle \triangle }
150
+
151
+ AOD is equal to that of
152
+
153
+
154
+
155
+ △
156
+
157
+
158
+ {\displaystyle \triangle }
159
+
160
+ BOC, and the product of the areas of
161
+
162
+
163
+
164
+ △
165
+
166
+
167
+ {\displaystyle \triangle }
168
+
169
+ AOD and
170
+
171
+
172
+
173
+ △
174
+
175
+
176
+ {\displaystyle \triangle }
177
+
178
+ BOC is equal to that of
179
+
180
+
181
+
182
+ △
183
+
184
+
185
+ {\displaystyle \triangle }
186
+
187
+ AOB and
188
+
189
+
190
+
191
+ △
192
+
193
+
194
+ {\displaystyle \triangle }
195
+
196
+ COD. The ratio of the areas of each pair of adjacent triangles is the same as that between the lengths of the parallel sides.[11]
197
+
198
+ Let the trapezoid have vertices A, B, C, and D in sequence and have parallel sides AB and DC. Let E be the intersection of the diagonals, and let F be on side DA and G be on side BC such that FEG is parallel to AB and CD. Then FG is the harmonic mean of AB and DC:[16]
199
+
200
+ The line that goes through both the intersection point of the extended nonparallel sides and the intersection point of the diagonals, bisects each base.[17]
201
+
202
+ The center of area (center of mass for a uniform lamina) lies along the line segment joining the midpoints of the parallel sides, at a perpendicular distance x from the longer side b given by[18]
203
+
204
+ The center of area divides this segment in the ratio (when taken from the short to the long side)[19]:p. 862
205
+
206
+ If the angle bisectors to angles A and B intersect at P, and the angle bisectors to angles C and D intersect at Q, then[17]
207
+
208
+ In architecture the word is used to refer to symmetrical doors, windows, and buildings built wider at the base, tapering toward the top, in Egyptian style. If these have straight sides and sharp angular corners, their shapes are usually isosceles trapezoids. This was the standard style for the doors and windows of the Inca.[20]
209
+
210
+ The crossed ladders problem is the problem of finding the distance between the parallel sides of a right trapezoid, given the diagonal lengths and the distance from the perpendicular leg to the diagonal intersection.
211
+
212
+ In morphology, taxonomy and other descriptive disciplines in which a term for such shapes is necessary, terms such as trapezoidal or trapeziform commonly are useful in descriptions of particular organs or forms.[21]
213
+
214
+ In computer engineering, specifically digital logic and computer architecture, trapezoids are typically utilized to symbolize multiplexors. Multiplexors are logic elements that select between multiple elements and produce a single output based on a select signal. Typical designs will employ trapezoids without specifically stating they are multiplexors as they are universally equivalent.
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1
+ The Twelve Labours of Heracles or Hercules (Greek: οἱ Ἡρακλέους ἆθλοι, hoi Hērakléous âthloi)[1][2] are a series of episodes concerning a penance carried out by Heracles, the greatest of the Greek heroes, whose name was later romanised as Hercules. They were accomplished at the service of King Eurystheus. The episodes were later connected by a continuous narrative. The establishment of a fixed cycle of twelve labours was attributed by the Greeks to an epic poem, now lost, written by Peisander, dated about 600 BC.[3] After Heracles killed his wife and children, he went to the oracle at Delphi. He prayed to the god Apollo for guidance. Heracles was told to serve the king of Mycenae, Eurystheus, for ten years. During this time, he is sent to perform a series of difficult feats, called labours.[4]
2
+
3
+ Driven mad by Hera (queen of the gods), Heracles slew his sons by his wife Megara.[5] After recovering his sanity, Heracles deeply regretted his actions; he was purified by King Thespius, then traveled to Delphi to inquire how he could atone for his actions. Pythia, the Oracle of Delphi, advised him to go to Tiryns and serve his cousin, King Eurystheus, for ten years, performing whatever labours Eurystheus might set him; in return, he would be rewarded with immortality. Heracles despaired at this, loathing to serve a man whom he knew to be far inferior to himself, yet fearing to oppose his father Zeus. Eventually, he placed himself at Eurystheus's disposal.
4
+
5
+ Eurystheus originally ordered Heracles to perform ten labours. Heracles accomplished these tasks, but Eurystheus refused to recognize two: the slaying of the Lernaean Hydra, as Heracles' nephew and charioteer Iolaus had helped him; and the cleansing of the Augeas, because Heracles accepted payment for the labour. Eurystheus set two more tasks (fetching the Golden Apples of Hesperides and capturing Cerberus), which Heracles also performed, bringing the total number of tasks to twelve.
6
+
7
+ As they survive, the labours of Heracles are not recounted in any single place, but must be reassembled from many sources. Ruck and Staples[6] assert that there is no one way to interpret the labours, but that six were located in the Peloponnese, culminating with the rededication of Olympia. Six others took the hero farther afield, to places that were, per Ruck, "all previously strongholds of Hera or the 'Goddess' and were Entrances to the Netherworld".[6] In each case, the pattern was the same: Heracles was sent to kill or subdue, or to fetch back for Eurystheus (as Hera's representative) a magical animal or plant.
8
+
9
+ A famous depiction of the labours in Greek sculpture is found on the metopes of the Temple of Zeus at Olympia, which date to the 450s BC.[citation needed]
10
+
11
+ In his labours, Heracles was sometimes accompanied by a male companion (an eromenos), according to Licymnius[citation needed] and others, such as Iolaus, his nephew. Although he was supposed to perform only ten labours, this assistance led to two labours being disqualified: Eurystheus refused to recognize slaying the Hydra, because Iolaus helped him, and the cleansing of the Augean stables, because Heracles was paid for his services and/or because the rivers did the work. Several of the labours involved the offspring (by various accounts) of Typhon and his mate Echidna, all overcome by Heracles.
12
+
13
+ A traditional order of the labours found in the Bibliotheca[7] by Pseudo-Apollodorus is:
14
+
15
+ Heracles wandered the area until he came to the town of Cleonae. There he met a boy who said that if Heracles slew the Nemean lion and returned alive within 30 days, the town would sacrifice a lion to Zeus, but if he did not return within 30 days or if he died, the boy would sacrifice himself to Zeus. Another version claims that he met Molorchos, a shepherd who had lost his son to the lion, saying that if he came back within 30 days, a ram would be sacrificed to Zeus. If he did not return within 30 days, it would be sacrificed to the dead Heracles as a mourning offering.
16
+
17
+ While searching for the lion, Heracles fletched some arrows to use against it, not knowing that its golden fur was impenetrable. When he found and shot the lion, firing at it with his bow, he discovered the fur's protective property as the arrow bounced harmlessly off the creature's thigh. After some time, Heracles made the lion return to his cave. The cave had two entrances, one of which Heracles blocked; he then entered the other. In those dark and close quarters, Heracles stunned the beast with his club and, using his immense strength, strangled it to death. During the fight the lion bit off one of his fingers.[8] Others say that he shot arrows at it, eventually shooting it in the unarmored mouth. After slaying the lion, he tried to skin it with a knife from his belt, but failed. He then tried sharpening the knife with a stone and even tried with the stone itself. Finally, Athena, noticing the hero's plight, told Heracles to use one of the lion's own claws to skin the pelt. Others say that Heracles' armor was, in fact, the hide of the Lion of Cithaeron.
18
+
19
+ When he returned on the 30th day carrying the carcass of the lion on his shoulders, King Eurystheus was amazed and terrified. Eurystheus forbade him ever again to enter the city; from then on he was to display the fruits of his labours outside the city gates. Eurystheus would then tell Heracles his tasks through a herald, not personally. Eurystheus even had a large bronze jar made for him in which to hide from Heracles if need be. Eurystheus then warned him that the tasks would become increasingly difficult.
20
+
21
+ Heracles' second labour was to slay the Lernaean Hydra, which Hera had raised just to slay Heracles. Upon reaching the swamp near Lake Lerna, where the Hydra dwelt, Heracles used a cloth to cover his mouth and nose to protect himself from the poisonous fumes. He fired flaming arrows into the Hydra's lair, the spring of Amymone, a deep cave that it only came out of to terrorize neighboring villages.[9] He then confronted the Hydra, wielding a harvesting sickle (according to some early vase-paintings), a sword or his famed club. Ruck and Staples (1994: 170) have pointed out that the chthonic creature's reaction was botanical: upon cutting off each of its heads he found that two grew back, an expression of the hopelessness of such a struggle for any but the hero. Additionally, one of the Hydra's heads - the middle one - was immortal.
22
+
23
+ The details of the struggle are explicit in the Bibliotheca (2.5.2): realizing that he could not defeat the Hydra in this way, Heracles called on his nephew Iolaus for help. His nephew then came upon the idea (possibly inspired by Athena) of using a firebrand to scorch the neck stumps after each decapitation. Heracles cut off each head and Iolaus cauterized the open stumps. Seeing that Heracles was winning the struggle, Hera sent a giant crab to distract him. He crushed it under his mighty foot. He cut off the Hydra's one immortal head with a golden sword given to him by Athena. Heracles placed it under a great rock on the sacred way between Lerna and Elaius (Kerenyi 1959:144), and dipped his arrows in the Hydra's poisonous blood, and so his second task was complete. The alternative version of this myth is that after cutting off one head, he then dipped his sword in it and used its venom to burn each head so it could not grow back. Hera, upset that Heracles had slain the beast she raised to kill him, placed it in the dark blue vault of the sky as the constellation Hydra. She then turned the crab into the constellation Cancer.
24
+
25
+ Later, Heracles used an arrow dipped in the Hydra's poisonous blood to kill the centaur Nessus; and Nessus's tainted blood was applied to the Tunic of Nessus, by which the centaur had his posthumous revenge. Both Strabo and Pausanias report that the stench of the river Anigrus in Elis, making all the fish of the river inedible, was reputed to be due to the Hydra's venom, washed from the arrows Heracles used on the centaur.[10]
26
+
27
+ Eurystheus and Hera were greatly angered that Heracles had survived the Nemean Lion and the Lernaean Hydra. For the third labour, they found a task which they thought would spell doom for the hero. It was not slaying a beast or monster, as it had already been established that Heracles could overcome even the most fearsome opponents. Instead, Eurystheus ordered him to capture the Ceryneian Hind, which was so fast that it could outrun an arrow.
28
+
29
+ After beginning the search, Heracles awoke from sleeping and saw the hind by the glint on its antlers. Heracles then chased the hind on foot for a full year through Greece, Thrace, Istria, and the land of the Hyperboreans. In some versions, he captured the hind while it slept, rendering it lame with a trap net. In other versions, he encountered Artemis in her temple; she told him to leave the hind and tell Eurystheus all that had happened, and his third labour would be considered to be completed. Yet another version claims that Heracles trapped the Hind with an arrow between its forelegs.
30
+
31
+ Eurystheus had given Heracles this task hoping to incite Artemis' anger at Heracles for his desecration of her sacred animal. As he was returning with the hind, Heracles encountered Artemis and her brother Apollo. He begged the goddess for forgiveness, explaining that he had to catch it as part of his penance, but he promised to return it. Artemis forgave him, foiling Eurystheus' plan to have her punish him.
32
+
33
+ Upon bringing the hind to Eurystheus, he was told that it was to become part of the King's menagerie. Heracles knew that he had to return the hind as he had promised, so he agreed to hand it over on the condition that Eurystheus himself come out and take it from him. The King came out, but the moment that Heracles let the hind go, it sprinted back to its mistress and Heracles left, saying that Eurystheus had not been quick enough.
34
+
35
+ Eurystheus was disappointed that Heracles had overcome yet another creature and was humiliated by the hind's escape, so he assigned Heracles another dangerous task. By some accounts, the fourth labour was to bring the fearsome Erymanthian Boar back to Eurystheus alive (there is no single definitive telling of the labours). On the way to Mount Erymanthos where the boar lived, Heracles visited Pholus ("caveman"), a kind and hospitable centaur and old friend. Heracles ate with Pholus in his cavern (though the centaur devoured his meat raw) and asked for wine. Pholus had only one jar of wine, a gift from Dionysus to all the centaurs on Mount Erymanthos. Heracles convinced him to open it, and the smell attracted the other centaurs. They did not understand that wine needs to be tempered with water, became drunk, and attacked Heracles. Heracles shot at them with his poisonous arrows, killing many, and the centaurs retreated all the way to Chiron's cave.
36
+
37
+ Pholus was curious why the arrows caused so much death. He picked one up but dropped it, and the arrow stabbed his hoof, poisoning him. One version states that a stray arrow hit Chiron as well. He was immortal, but he still felt the pain. Chiron's pain was so great that he volunteered to give up his immortality and take the place of Prometheus, who had been chained to the top of a mountain to have his liver eaten daily by an eagle. Prometheus' torturer, the eagle, continued its torture on Chiron, so Heracles shot it dead with an arrow. It is generally accepted that the tale was meant to show Heracles as being the recipient of Chiron's surrendered immortality. However, this tale contradicts the fact that Chiron later taught Achilles. The tale of the centaurs sometimes appears in other parts of the twelve labours, as does the freeing of Prometheus.
38
+
39
+ Heracles had visited Chiron to gain advice on how to catch the boar, and Chiron had told him to drive it into thick snow, which sets this labour in mid-winter. Heracles caught the boar, bound it, and carried it back to Eurystheus, who was frightened of it and ducked down in his half-buried storage pithos, begging Heracles to get rid of the beast.
40
+
41
+ The fifth labour was to clean the stables of King Augeas. This assignment was intended to be both humiliating (rather than impressive, as the previous labours had been) and impossible, since the livestock were divinely healthy (and immortal) and therefore produced an enormous quantity of dung. The Augean Stables (/ɔːˈdʒiːən/) had not been cleaned in over 30 years, and over 1,000 cattle lived there. However, Heracles succeeded by re-routing the rivers Alpheus and Peneus to wash out the filth.
42
+
43
+ Before starting on the task, Heracles had asked Augeas for one-tenth of the cattle if he finished the task in one day, and Augeas agreed. But afterwards Augeas refused to honour the agreement on the grounds that Heracles had been ordered to carry out the task by Eurystheus anyway. Heracles claimed his reward in court, and was supported by Augeas' son Phyleus. Augeas banished them both before the court had ruled. Heracles returned, slew Augeas, and gave his kingdom to Phyleus. Heracles then founded the Olympic Games.
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+ The success of this labour was ultimately discounted as the rushing waters had done the work of cleaning the stables and because Heracles was paid for doing the labour.
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+ Eurystheus said that Heracles still had seven labours to perform.[11]
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+ The sixth labour was to defeat the Stymphalian birds, man-eating birds with beaks made of bronze and sharp metallic feathers they could launch at their victims. They were sacred to Ares, the god of war. Furthermore, their dung was highly toxic. They had migrated to Lake Stymphalia in Arcadia, where they bred quickly and took over the countryside, destroying local crops, fruit trees, and townspeople. Heracles could not go too far into the swamp, for it would not support his weight. Athena, noticing the hero's plight, gave Heracles a rattle which Hephaestus had made especially for the occasion. Heracles shook the rattle and frightened the birds into the air. Heracles then shot many of them with his arrows. The rest flew far away, never to return. The Argonauts would later encounter them.
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+ The seventh labour was to capture the Cretan Bull, father of the Minotaur. Heracles sailed to Crete, where King Minos gave Heracles permission to take the bull away and even offered him assistance (which Heracles declined, plausibly because he did not want the labour to be discounted as before).[12] The bull had been wreaking havoc on Crete by uprooting crops and leveling orchard walls. Heracles sneaked up behind the bull and then used his hands to throttle it (stopping before it was killed), and then shipped it back to Tiryns. Eurystheus, who hid in his pithos at first sight of the creature, wanted to sacrifice the bull to Hera, who hated Heracles. She refused the sacrifice because it reflected glory on Heracles. The bull was released and wandered into Marathon, becoming known as the Marathonian Bull.[12] Theseus would later sacrifice the bull to Athena and/or Apollo.
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+ As the eighth of his Twelve Labours, also categorised as the second of the Non-Peloponneisan labours,[13] Heracles was sent by King Eurystheus to steal the Mares from Diomedes. The mares’ madness was attributed to their unnatural diet which consisted of the flesh[14] of unsuspecting guests or strangers to the island.[15] Some versions of the myth say that the mares also expelled fire when they breathed.[16] The Mares, which were the terror of Thrace, were kept tethered by iron chains to a bronze manger in the now vanished city of Tirida[17] and were named Podargos (the swift), Lampon (the shining), Xanthos (the yellow) and Deinos (or Deinus, the terrible).[18] Although very similar, there are slight variances in the exact details regarding the mares’ capture.
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+ In one version, Heracles brought a number of volunteers to help him capture the giant horses.[17] After overpowering Diomedes’ men, Heracles broke the chains that tethered the horses and drove the mares down to sea. Unaware that the mares were man-eating and uncontrollable, Heracles left them in the charge of his favored companion, Abderus, while he left to fight Diomedes. Upon his return, Heracles found that the boy was eaten. As revenge, Heracles fed Diomedes to his own horses and then founded Abdera next to the boy's tomb.[15]
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+ In another version, Heracles, who was visiting the island, stayed awake so that he didn't have his throat cut by Diomedes in the night, and cut the chains binding the horses once everyone was asleep. Having scared the horses onto the high ground of a knoll, Heracles quickly dug a trench through the peninsula, filling it with water and thus flooding the low-lying plain. When Diomedes and his men turned to flee, Heracles killed them with an axe (or a club[17]), and fed Diomedes’ body to the horses to calm them.
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+ In yet another version, Heracles first captured Diomedes and fed him to the mares before releasing them. Only after realizing that their King was dead did his men, the Bistonians,[15][17] attack Heracles. Upon seeing the mares charging at them, led in a chariot by Abderus, the Bistonians turned and fled.
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+ All versions have eating human flesh make the horses calmer, giving Heracles the opportunity to bind their mouths shut, and easily take them back to King Eurystheus, who dedicated the horses to Hera.[19] In some versions, they were allowed to roam freely around Argos, having become permanently calm, but in others, Eurystheus ordered the horses taken to Olympus to be sacrificed to Zeus, but Zeus refused them, and sent wolves, lions, and bears to kill them.[20] Roger Lancelyn Green states in his Tales of the Greek Heroes that the mares’ descendants were used in the Trojan War, and survived even to the time of Alexander the Great.[17][21] After the incident, Eurystheus sent Heracles to bring back Hippolyta's Girdle.
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+ Eurystheus' daughter Admete wanted the Belt of Hippolyta, queen of the Amazons, a gift from her father Ares. To please his daughter, Eurystheus ordered Heracles to retrieve the belt as his ninth labour.
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+ Taking a band of friends with him, Heracles set sail, stopping at the island of Paros, which was inhabited by some sons of Minos. The sons killed two of Heracles' companions, an act which set Heracles on a rampage. He killed two of the sons of Minos and threatened the other inhabitants until he was offered two men to replace his fallen companions. Heracles agreed and took two of Minos' grandsons, Alcaeus and Sthenelus. They continued their voyage and landed at the court of Lycus, whom Heracles defended in a battle against King Mygdon of Bebryces. After killing King Mygdon, Heracles gave much of the land to his friend Lycus. Lycus called the land Heraclea. The crew then set off for Themiscyra, where Hippolyta lived.
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+ All would have gone well for Heracles had it not been for Hera. Hippolyta, impressed with Heracles and his exploits, agreed to give him the belt and would have done so had Hera not disguised herself and walked among the Amazons sowing seeds of distrust. She claimed the strangers were plotting to carry off the queen of the Amazons. Alarmed, the women set off on horseback to confront Heracles. When Heracles saw them, he thought Hippolyta had been plotting such treachery all along and had never meant to hand over the belt, so he killed her, took the belt and returned to Eurystheus.
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+ The tenth labour was to obtain the Cattle of the three-bodied giant Geryon. In the fullest account in the Bibliotheca of Pseudo-Apollodorus,[22] Heracles had to go to the island of Erytheia in the far west (sometimes identified with the Hesperides, or with the island which forms the city of Cádiz) to get the cattle. On the way there, he crossed the Libyan desert[23] and became so frustrated at the heat that he shot an arrow at the Sun. The sun-god Helios "in admiration of his courage" gave Heracles the golden cup Helios used to sail across the sea from west to east each night. Heracles rode the cup to Erytheia; Heracles in the cup was a favorite motif on black-figure pottery.[citation needed] Such a magical conveyance undercuts any literal geography for Erytheia, the "red island" of the sunset.
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+ When Heracles landed at Erytheia, he was confronted by the two-headed dog Orthrus. With one blow from his olive-wood club, Heracles killed Orthrus. Eurytion the herdsman came to assist Orthrus, but Heracles dealt with him the same way.
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+ On hearing the commotion, Geryon sprang into action, carrying three shields and three spears, and wearing three helmets. He attacked Heracles at the River Anthemus, but was slain by one of Heracles' poisoned arrows. Heracles shot so forcefully that the arrow pierced Geryon's forehead, "and Geryon bent his neck over to one side, like a poppy that spoils its delicate shapes, shedding its petals all at once."[24]
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+ Heracles then had to herd the cattle back to Eurystheus. In Roman versions of the narrative, Heracles drove the cattle over the Aventine Hill on the future site of Rome. The giant Cacus, who lived there, stole some of the cattle as Heracles slept, making the cattle walk backwards so that they left no trail, a repetition of the trick of the young Hermes. According to some versions, Heracles drove his remaining cattle past the cave, where Cacus had hidden the stolen animals, and they began calling out to each other. In other versions, Cacus' sister Caca told Heracles where he was. Heracles then killed Cacus, and set up an altar on the spot, later the site of Rome's Forum Boarium (the cattle market).
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+ To annoy Heracles, Hera sent a gadfly to bite the cattle, irritate them, and scatter them. Within a year, Heracles retrieved them. Hera then sent a flood which raised the level of a river so much that Heracles could not cross with the cattle. He piled stones into the river to make the water shallower. When he finally reached the court of Eurystheus, the cattle were sacrificed to Hera.
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+ After Heracles completed the first ten labours, Eurystheus gave him two more, claiming that slaying the Hydra did not count (because Iolaus helped Heracles), neither did cleaning the Augean Stables (either because he was paid for the job or because the rivers did the work).
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+ The first additional labour was to steal three of the golden apples from the garden of the Hesperides. Heracles first caught the Old Man of the Sea, the shapeshifting sea god,[25] to learn where the Garden of the Hesperides was located.[26]
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+ In some variations, Heracles, either at the start or at the end of this task, meets Antaeus, who was invincible as long as he touched his mother, Gaia, the Earth. Heracles killed Antaeus by holding him aloft and crushing him in a bear hug.[27]
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+ Herodotus claims that Heracles stopped in Egypt, where King Busiris decided to make him the yearly sacrifice, but Heracles burst out of his chains.
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+ Heracles finally made his way to the garden of the Hesperides, where he encountered Atlas holding up the heavens on his shoulders. Heracles persuaded Atlas to get the three golden Apples for him by offering to hold up the heavens in his place for a little while. Atlas could get the apples because, in this version, he was the father or otherwise related to the Hesperides. This would have made the labour – like the Hydra and the Augean stables – void because Heracles had received help. When Atlas returned, he decided that he did not want to take the heavens back, and instead offered to deliver the apples himself, but Heracles tricked him by agreeing to remain in place of Atlas on the condition that Atlas relieve him temporarily while Heracles adjusted his cloak. Atlas agreed, but Heracles reneged and walked away with the apples. According to an alternative version, Heracles slew Ladon, the dragon who guarded the apples instead. Eurystheus was furious that Heracles had accomplished something that Eurystheus thought could not possibly be done.
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+ The twelfth and final labour was the capture of Cerberus, the three-headed, dragon-tailed dog that was the guardian of the gates of the Underworld. To prepare for his descent into the Underworld, Heracles went to Eleusis (or Athens) to be initiated in the Eleusinian Mysteries. He entered the Underworld, and Hermes and Athena were his guides.
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+ While in the Underworld, Heracles met Theseus and Pirithous. The two companions had been imprisoned by Hades for attempting to kidnap Persephone. One tradition tells of snakes coiling around their legs, then turning into stone; another that Hades feigned hospitality and prepared a feast inviting them to sit. They unknowingly sat in chairs of forgetfulness and were permanently ensnared. When Heracles had pulled Theseus first from his chair, some of his thigh stuck to it (this explains the supposedly lean thighs of Athenians), but the Earth shook at the attempt to liberate Pirithous, whose desire to have the goddess for himself was so insulting he was doomed to stay behind.
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+ Heracles found Hades and asked permission to bring Cerberus to the surface, which Hades agreed to if Heracles could subdue the beast without using weapons. Heracles overpowered Cerberus with his bare hands and slung the beast over his back. He carried Cerberus out of the Underworld through a cavern entrance in the Peloponnese and brought it to Eurystheus, who again fled into his pithos. Eurystheus begged Heracles to return Cerberus to the Underworld, offering in return to release him from any further labours when Cerberus disappeared back to his master.
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+ After completing the Twelve Labours, one tradition says Heracles joined Jason and the Argonauts in their quest for the Golden Fleece. However, Herodorus (c. 400 BC) disputed this and denied Heracles ever sailed with the Argonauts. A separate tradition (e.g. Argonautica) has Heracles accompany the Argonauts, but he did not travel with them as far as Colchis.
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+ Some ancient Greeks found allegorical meanings of a moral, psychological or philosophical nature in the Labours of Heracles. This trend became more prominent in the Renaissance.[28] For example, Heraclitus the Grammarian wrote in his Homeric Problems:
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+ I turn to Heracles. We must not suppose he attained such power in those days as a result of his physical strength. Rather, he was a man of intellect, an initiate in heavenly wisdom, who, as it were, shed light on philosophy, which had been hidden in deep darkness. The most authoritative of the Stoics agree with this account.... The (Erymanthian) boar which he overcame is the common incontinence of men; the (Nemean) lion is the indiscriminate rush towards improper goals; in the same way, by fettering irrational passions he gave rise to the belief that he had fettered the violent (Cretan) bull. He banished cowardice also from the world, in the shape of the hind of Ceryneia. There was another "labor" too, not properly so called, in which he cleared out the mass of dung (from the Augean stables) — in other words, the foulness that disfigures humanity. The (Stymphalian) birds he scattered are the windy hopes that feed our lives; the many-headed hydra that he burned, as it were, with the fires of exhortation, is pleasure, which begins to grow again as soon as it is cut out.
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+ Color (American English), or colour (Commonwealth English), is the characteristic of visual perception described through color categories, with names such as red, orange, yellow, green, blue, or purple. This perception of color derives from the stimulation of photoreceptor cells (in particular cone cells in the human eye and other vertebrate eyes) by electromagnetic radiation (in the visible spectrum in the case of humans). Color categories and physical specifications of color are associated with objects through the wavelengths of the light that is reflected from them and their intensities. This reflection is governed by the object's physical properties such as light absorption, emission spectra, etc.
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+ By defining a color space, colors can be identified numerically by coordinates, which in 1931 were also named in global agreement with internationally agreed color names like mentioned above (red, orange, etc.) by the International Commission on Illumination. The RGB color space for instance is a color space corresponding to human trichromacy and to the three cone cell types that respond to three bands of light: long wavelengths, peaking near 564–580 nm (red); medium-wavelength, peaking near 534–545 nm (green); and short-wavelength light, near 420–440 nm (blue).[1][2] There may also be more than three color dimensions in other color spaces, such as in the CMYK color model, wherein one of the dimensions relates to a color's colorfulness).
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+ The photo-receptivity of the "eyes" of other species also varies considerably from that of humans and so results in correspondingly different color perceptions that cannot readily be compared to one another. Honey bees and bumblebees for instance have trichromatic color vision sensitive to ultraviolet but is insensitive to red. Papilio butterflies possess six types of photoreceptors and may have pentachromatic vision.[3] The most complex color vision system in the animal kingdom has been found in stomatopods (such as the mantis shrimp) with up to 12 spectral receptor types thought to work as multiple dichromatic units.[4]
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+ The science of color is sometimes called chromatics, colorimetry, or simply color science. It includes the study of the perception of color by the human eye and brain, the origin of color in materials, color theory in art, and the physics of electromagnetic radiation in the visible range (that is, what is commonly referred to simply as light).
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+ Electromagnetic radiation is characterized by its wavelength (or frequency) and its intensity. When the wavelength is within the visible spectrum (the range of wavelengths humans can perceive, approximately from 390 nm to 700 nm), it is known as "visible light".
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+ Most light sources emit light at many different wavelengths; a source's spectrum is a distribution giving its intensity at each wavelength. Although the spectrum of light arriving at the eye from a given direction determines the color sensation in that direction, there are many more possible spectral combinations than color sensations. In fact, one may formally define a color as a class of spectra that give rise to the same color sensation, although such classes would vary widely among different species, and to a lesser extent among individuals within the same species. In each such class the members are called metamers of the color in question. This effect can be visualized by comparing the light sources' spectral power distributions and the resulting colors.
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+ The familiar colors of the rainbow in the spectrum—named using the Latin word for appearance or apparition by Isaac Newton in 1671—include all those colors that can be produced by visible light of a single wavelength only, the pure spectral or monochromatic colors. The table at right shows approximate frequencies (in terahertz) and wavelengths (in nanometers) for various pure spectral colors. The wavelengths listed are as measured in air or vacuum (see refractive index).
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+ The color table should not be interpreted as a definitive list—the pure spectral colors form a continuous spectrum, and how it is divided into distinct colors linguistically is a matter of culture and historical contingency (although people everywhere have been shown to perceive colors in the same way[6]). A common list identifies six main bands: red, orange, yellow, green, blue, and violet. Newton's conception included a seventh color, indigo, between blue and violet. It is possible that what Newton referred to as blue is nearer to what today is known as cyan, and that indigo was simply the dark blue of the indigo dye that was being imported at the time.[7]
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+ The intensity of a spectral color, relative to the context in which it is viewed, may alter its perception considerably; for example, a low-intensity orange-yellow is brown, and a low-intensity yellow-green is olive green.
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+ The color of an object depends on both the physics of the object in its environment and the characteristics of the perceiving eye and brain. Physically, objects can be said to have the color of the light leaving their surfaces, which normally depends on the spectrum of the incident illumination and the reflectance properties of the surface, as well as potentially on the angles of illumination and viewing. Some objects not only reflect light, but also transmit light or emit light themselves, which also contributes to the color. A viewer's perception of the object's color depends not only on the spectrum of the light leaving its surface, but also on a host of contextual cues, so that color differences between objects can be discerned mostly independent of the lighting spectrum, viewing angle, etc. This effect is known as color constancy.
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+ Some generalizations of the physics can be drawn, neglecting perceptual effects for now:
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+ To summarize, the color of an object is a complex result of its surface properties, its transmission properties, and its emission properties, all of which contribute to the mix of wavelengths in the light leaving the surface of the object. The perceived color is then further conditioned by the nature of the ambient illumination, and by the color properties of other objects nearby, and via other characteristics of the perceiving eye and brain.
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+ Although Aristotle and other ancient scientists had already written on the nature of light and color vision, it was not until Newton that light was identified as the source of the color sensation. In 1810, Goethe published his comprehensive Theory of Colors in which he ascribed physiological effects to color that are now understood as psychological.
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+ In 1801 Thomas Young proposed his trichromatic theory, based on the observation that any color could be matched with a combination of three lights. This theory was later refined by James Clerk Maxwell and Hermann von Helmholtz. As Helmholtz puts it, "the principles of Newton's law of mixture were experimentally confirmed by Maxwell in 1856. Young's theory of color sensations, like so much else that this marvelous investigator achieved in advance of his time, remained unnoticed until Maxwell directed attention to it."[10]
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+ At the same time as Helmholtz, Ewald Hering developed the opponent process theory of color, noting that color blindness and afterimages typically come in opponent pairs (red-green, blue-orange, yellow-violet, and black-white). Ultimately these two theories were synthesized in 1957 by Hurvich and Jameson, who showed that retinal processing corresponds to the trichromatic theory, while processing at the level of the lateral geniculate nucleus corresponds to the opponent theory.[11]
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+ In 1931, an international group of experts known as the Commission internationale de l'éclairage (CIE) developed a mathematical color model, which mapped out the space of observable colors and assigned a set of three numbers to each.
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+ The ability of the human eye to distinguish colors is based upon the varying sensitivity of different cells in the retina to light of different wavelengths. Humans are trichromatic—the retina contains three types of color receptor cells, or cones. One type, relatively distinct from the other two, is most responsive to light that is perceived as blue or blue-violet, with wavelengths around 450 nm; cones of this type are sometimes called short-wavelength cones or S cones (or misleadingly, blue cones). The other two types are closely related genetically and chemically: middle-wavelength cones, M cones, or green cones are most sensitive to light perceived as green, with wavelengths around 540 nm, while the long-wavelength cones, L cones, or red cones, are most sensitive to light that is perceived as greenish yellow, with wavelengths around 570 nm.
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+ Light, no matter how complex its composition of wavelengths, is reduced to three color components by the eye. Each cone type adheres to the principle of univariance, which is that each cone's output is determined by the amount of light that falls on it over all wavelengths. For each location in the visual field, the three types of cones yield three signals based on the extent to which each is stimulated. These amounts of stimulation are sometimes called tristimulus values.
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+ The response curve as a function of wavelength varies for each type of cone. Because the curves overlap, some tristimulus values do not occur for any incoming light combination. For example, it is not possible to stimulate only the mid-wavelength (so-called "green") cones; the other cones will inevitably be stimulated to some degree at the same time. The set of all possible tristimulus values determines the human color space. It has been estimated that humans can distinguish roughly 10 million different colors.[9]
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+ The other type of light-sensitive cell in the eye, the rod, has a different response curve. In normal situations, when light is bright enough to strongly stimulate the cones, rods play virtually no role in vision at all.[12] On the other hand, in dim light, the cones are understimulated leaving only the signal from the rods, resulting in a colorless response. (Furthermore, the rods are barely sensitive to light in the "red" range.) In certain conditions of intermediate illumination, the rod response and a weak cone response can together result in color discriminations not accounted for by cone responses alone. These effects, combined, are summarized also in the Kruithof curve, that describes the change of color perception and pleasingness of light as function of temperature and intensity.
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+ While the mechanisms of color vision at the level of the retina are well-described in terms of tristimulus values, color processing after that point is organized differently. A dominant theory of color vision proposes that color information is transmitted out of the eye by three opponent processes, or opponent channels, each constructed from the raw output of the cones: a red–green channel, a blue–yellow channel, and a black–white "luminance" channel. This theory has been supported by neurobiology, and accounts for the structure of our subjective color experience. Specifically, it explains why humans cannot perceive a "reddish green" or "yellowish blue", and it predicts the color wheel: it is the collection of colors for which at least one of the two color channels measures a value at one of its extremes.
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+ The exact nature of color perception beyond the processing already described, and indeed the status of color as a feature of the perceived world or rather as a feature of our perception of the world—a type of qualia—is a matter of complex and continuing philosophical dispute.
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+ If one or more types of a person's color-sensing cones are missing or less responsive than normal to incoming light, that person can distinguish fewer colors and is said to be color deficient or color blind (though this latter term can be misleading; almost all color deficient individuals can distinguish at least some colors). Some kinds of color deficiency are caused by anomalies in the number or nature of cones in the retina. Others (like central or cortical achromatopsia) are caused by neural anomalies in those parts of the brain where visual processing takes place.
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+ While most humans are trichromatic (having three types of color receptors), many animals, known as tetrachromats, have four types. These include some species of spiders, most marsupials, birds, reptiles, and many species of fish. Other species are sensitive to only two axes of color or do not perceive color at all; these are called dichromats and monochromats respectively. A distinction is made between retinal tetrachromacy (having four pigments in cone cells in the retina, compared to three in trichromats) and functional tetrachromacy (having the ability to make enhanced color discriminations based on that retinal difference). As many as half of all women are retinal tetrachromats.[13]:p.256 The phenomenon arises when an individual receives two slightly different copies of the gene for either the medium- or long-wavelength cones, which are carried on the X chromosome. To have two different genes, a person must have two X chromosomes, which is why the phenomenon only occurs in women.[13] There is one scholarly report that confirms the existence of a functional tetrachromat.[14]
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+ In certain forms of synesthesia/ideasthesia, perceiving letters and numbers (grapheme–color synesthesia) or hearing musical sounds (music–color synesthesia) will lead to the unusual additional experiences of seeing colors. Behavioral and functional neuroimaging experiments have demonstrated that these color experiences lead to changes in behavioral tasks and lead to increased activation of brain regions involved in color perception, thus demonstrating their reality, and similarity to real color percepts, albeit evoked through a non-standard route.
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+ After exposure to strong light in their sensitivity range, photoreceptors of a given type become desensitized. For a few seconds after the light ceases, they will continue to signal less strongly than they otherwise would. Colors observed during that period will appear to lack the color component detected by the desensitized photoreceptors. This effect is responsible for the phenomenon of afterimages, in which the eye may continue to see a bright figure after looking away from it, but in a complementary color.
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+ Afterimage effects have also been utilized by artists, including Vincent van Gogh.
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+ When an artist uses a limited color palette, the eye tends to compensate by seeing any gray or neutral color as the color which is missing from the color wheel. For example, in a limited palette consisting of red, yellow, black, and white, a mixture of yellow and black will appear as a variety of green, a mixture of red and black will appear as a variety of purple, and pure gray will appear bluish.[15]
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+ The trichromatic theory is strictly true when the visual system is in a fixed state of adaptation. In reality, the visual system is constantly adapting to changes in the environment and compares the various colors in a scene to reduce the effects of the illumination. If a scene is illuminated with one light, and then with another, as long as the difference between the light sources stays within a reasonable range, the colors in the scene appear relatively constant to us. This was studied by Edwin Land in the 1970s and led to his retinex theory of color constancy.
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+ Both phenomena are readily explained and mathematically modeled with modern theories of chromatic adaptation and color appearance (e.g. CIECAM02, iCAM).[16] There is no need to dismiss the trichromatic theory of vision, but rather it can be enhanced with an understanding of how the visual system adapts to changes in the viewing environment.
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+ Colors vary in several different ways, including hue (shades of red, orange, yellow, green, blue, and violet), saturation, brightness, and gloss. Some color words are derived from the name of an object of that color, such as "orange" or "salmon", while others are abstract, like "red".
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+ In the 1969 study Basic Color Terms: Their Universality and Evolution, Brent Berlin and Paul Kay describe a pattern in naming "basic" colors (like "red" but not "red-orange" or "dark red" or "blood red", which are "shades" of red). All languages that have two "basic" color names distinguish dark/cool colors from bright/warm colors. The next colors to be distinguished are usually red and then yellow or green. All languages with six "basic" colors include black, white, red, green, blue, and yellow. The pattern holds up to a set of twelve: black, gray, white, pink, red, orange, yellow, green, blue, purple, brown, and azure (distinct from blue in Russian and Italian, but not English).
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+ Colors, their meanings and associations can play major role in works of art, including literature.[17]
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+ Individual colors have a variety of cultural associations such as national colors (in general described in individual color articles and color symbolism). The field of color psychology attempts to identify the effects of color on human emotion and activity. Chromotherapy is a form of alternative medicine attributed to various Eastern traditions. Colors have different associations in different countries and cultures.[18]
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+ Different colors have been demonstrated to have effects on cognition. For example, researchers at the University of Linz in Austria demonstrated that the color red significantly decreases cognitive functioning in men.[19]
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+ Most light sources are mixtures of various wavelengths of light. Many such sources can still effectively produce a spectral color, as the eye cannot distinguish them from single-wavelength sources. For example, most computer displays reproduce the spectral color orange as a combination of red and green light; it appears orange because the red and green are mixed in the right proportions to allow the eye's cones to respond the way they do to the spectral color orange.
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+ A useful concept in understanding the perceived color of a non-monochromatic light source is the dominant wavelength, which identifies the single wavelength of light that produces a sensation most similar to the light source. Dominant wavelength is roughly akin to hue.
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+ There are many color perceptions that by definition cannot be pure spectral colors due to desaturation or because they are purples (mixtures of red and violet light, from opposite ends of the spectrum). Some examples of necessarily non-spectral colors are the achromatic colors (black, gray, and white) and colors such as pink, tan, and magenta.
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+ Two different light spectra that have the same effect on the three color receptors in the human eye will be perceived as the same color. They are metamers of that color. This is exemplified by the white light emitted by fluorescent lamps, which typically has a spectrum of a few narrow bands, while daylight has a continuous spectrum. The human eye cannot tell the difference between such light spectra just by looking into the light source, although reflected colors from objects can look different. (This is often exploited; for example, to make fruit or tomatoes look more intensely red.)
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+ Similarly, most human color perceptions can be generated by a mixture of three colors called primaries. This is used to reproduce color scenes in photography, printing, television, and other media. There are a number of methods or color spaces for specifying a color in terms of three particular primary colors. Each method has its advantages and disadvantages depending on the particular application.
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+ No mixture of colors, however, can produce a response truly identical to that of a spectral color, although one can get close, especially for the longer wavelengths, where the CIE 1931 color space chromaticity diagram has a nearly straight edge. For example, mixing green light (530 nm) and blue light (460 nm) produces cyan light that is slightly desaturated, because response of the red color receptor would be greater to the green and blue light in the mixture than it would be to a pure cyan light at 485 nm that has the same intensity as the mixture of blue and green.
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+ Because of this, and because the primaries in color printing systems generally are not pure themselves, the colors reproduced are never perfectly saturated spectral colors, and so spectral colors cannot be matched exactly. However, natural scenes rarely contain fully saturated colors, thus such scenes can usually be approximated well by these systems. The range of colors that can be reproduced with a given color reproduction system is called the gamut. The CIE chromaticity diagram can be used to describe the gamut.
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+ Another problem with color reproduction systems is connected with the acquisition devices, like cameras or scanners. The characteristics of the color sensors in the devices are often very far from the characteristics of the receptors in the human eye. In effect, acquisition of colors can be relatively poor if they have special, often very "jagged", spectra caused for example by unusual lighting of the photographed scene.
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+ A color reproduction system "tuned" to a human with normal color vision may give very inaccurate results for other observers.
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+ The different color response of different devices can be problematic if not properly managed. For color information stored and transferred in digital form, color management techniques, such as those based on ICC profiles, can help to avoid distortions of the reproduced colors. Color management does not circumvent the gamut limitations of particular output devices, but can assist in finding good mapping of input colors into the gamut that can be reproduced.
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+ Additive color is light created by mixing together light of two or more different colors. Red, green, and blue are the additive primary colors normally used in additive color systems such as projectors and computer terminals.
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+ Subtractive coloring uses dyes, inks, pigments, or filters to absorb some wavelengths of light and not others. The color that a surface displays comes from the parts of the visible spectrum that are not absorbed and therefore remain visible. Without pigments or dye, fabric fibers, paint base and paper are usually made of particles that scatter white light (all colors) well in all directions. When a pigment or ink is added, wavelengths are absorbed or "subtracted" from white light, so light of another color reaches the eye.
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+ If the light is not a pure white source (the case of nearly all forms of artificial lighting), the resulting spectrum will appear a slightly different color. Red paint, viewed under blue light, may appear black. Red paint is red because it scatters only the red components of the spectrum. If red paint is illuminated by blue light, it will be absorbed by the red paint, creating the appearance of a black object.
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+ Structural colors are colors caused by interference effects rather than by pigments. Color effects are produced when a material is scored with fine parallel lines, formed of one or more parallel thin layers, or otherwise composed of microstructures on the scale of the color's wavelength. If the microstructures are spaced randomly, light of shorter wavelengths will be scattered preferentially to produce Tyndall effect colors: the blue of the sky (Rayleigh scattering, caused by structures much smaller than the wavelength of light, in this case air molecules), the luster of opals, and the blue of human irises. If the microstructures are aligned in arrays, for example the array of pits in a CD, they behave as a diffraction grating: the grating reflects different wavelengths in different directions due to interference phenomena, separating mixed "white" light into light of different wavelengths. If the structure is one or more thin layers then it will reflect some wavelengths and transmit others, depending on the layers' thickness.
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+ Structural color is studied in the field of thin-film optics. The most ordered or the most changeable structural colors are iridescent. Structural color is responsible for the blues and greens of the feathers of many birds (the blue jay, for example), as well as certain butterfly wings and beetle shells. Variations in the pattern's spacing often give rise to an iridescent effect, as seen in peacock feathers, soap bubbles, films of oil, and mother of pearl, because the reflected color depends upon the viewing angle. Numerous scientists have carried out research in butterfly wings and beetle shells, including Isaac Newton and Robert Hooke. Since 1942, electron micrography has been used, advancing the development of products that exploit structural color, such as "photonic" cosmetics.[20]
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+ Basketball is a team sport in which two teams, most commonly of five players each, opposing one another on a rectangular court, compete with the primary objective of shooting a basketball (approximately 9.4 inches (24 cm) in diameter) through the defender's hoop (a basket 18 inches (46 cm) in diameter mounted 10 feet (3.048 m) high to a backboard at each end of the court) while preventing the opposing team from shooting through their own hoop. A field goal is worth two points, unless made from behind the three-point line, when it is worth three. After a foul, timed play stops and the player fouled or designated to shoot a technical foul is given one or more one-point free throws. The team with the most points at the end of the game wins, but if regulation play expires with the score tied, an additional period of play (overtime) is mandated.
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+ Players advance the ball by bouncing it while walking or running (dribbling) or by passing it to a teammate, both of which require considerable skill. On offense, players may use a variety of shots—the lay-up, the jump shot, or a dunk; on defense, they may steal the ball from a dribbler, intercept passes, or block shots; either offense or defense may collect a rebound, that is, a missed shot that bounces from rim or backboard. It is a violation to lift or drag one's pivot foot without dribbling the ball, to carry it, or to hold the ball with both hands then resume dribbling.
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+ The five players on each side fall into five playing positions. The tallest player is usually the center, the second tallest and strongest is the power forward, a slightly shorter but more agile player is the small forward, and the shortest players or the best ball handlers are the shooting guard and the point guard, who implements the coach's game plan by managing the execution of offensive and defensive plays (player positioning). Informally, players may play three-on-three, two-on-two, and one-on-one.
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+ Invented in 1891 by Canadian-American gym teacher James Naismith in Springfield, Massachusetts, United States, basketball has evolved to become one of the world's most popular and widely viewed sports.[1] The National Basketball Association (NBA) is the most significant professional basketball league in the world in terms of popularity, salaries, talent, and level of competition.[2][3] Outside North America, the top clubs from national leagues qualify to continental championships such as the EuroLeague and the Basketball Champions League Americas. The FIBA Basketball World Cup and Men's Olympic Basketball Tournament are the major international events of the sport and attract top national teams from around the world. Each continent hosts regional competitions for national teams, like EuroBasket and FIBA AmeriCup.
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+ The FIBA Women's Basketball World Cup and Women's Olympic Basketball Tournament feature top national teams from continental championships. The main North American league is the WNBA (NCAA Women's Division I Basketball Championship is also popular), whereas strongest European clubs participate in the EuroLeague Women.
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+ In early December 1891, Canadian James Naismith,[4] a physical education professor and instructor at the International Young Men's Christian Association Training School[5] (YMCA) (today, Springfield College) in Springfield, Massachusetts, was trying to keep his gym class active on a rainy day. He sought a vigorous indoor game to keep his students occupied and at proper levels of fitness during the long New England winters. After rejecting other ideas as either too rough or poorly suited to walled-in gymnasiums, he wrote the basic rules and nailed a peach basket onto an elevated track. In contrast with modern basketball nets, this peach basket retained its bottom, and balls had to be retrieved manually after each "basket" or point scored; this proved inefficient, however, so the bottom of the basket was removed, allowing the balls to be poked out with a long dowel each time.
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+ Basketball was originally played with a soccer ball. These round balls from "association football" were made, at the time, with a set of laces to close off the hole needed for inserting the inflatable bladder after the other sewn-together segments of the ball's cover had been flipped outside-in.[6][7] These laces could cause bounce passes and dribbling to be unpredictable.[8] Eventually a lace-free ball construction method was invented, and this change to the game was endorsed by Naismith. (Whereas in American football, the lace construction proved to be advantageous for gripping and remains to this day.) The first balls made specifically for basketball were brown, and it was only in the late 1950s that Tony Hinkle, searching for a ball that would be more visible to players and spectators alike, introduced the orange ball that is now in common use. Dribbling was not part of the original game except for the "bounce pass" to teammates. Passing the ball was the primary means of ball movement. Dribbling was eventually introduced but limited by the asymmetric shape of early balls.[dubious – discuss] Dribbling was common by 1896, with a rule against the double dribble by 1898.[9]
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+ The peach baskets were used until 1906 when they were finally replaced by metal hoops with backboards. A further change was soon made, so the ball merely passed through. Whenever a person got the ball in the basket, his team would gain a point. Whichever team got the most points won the game.[10] The baskets were originally nailed to the mezzanine balcony of the playing court, but this proved impractical when spectators in the balcony began to interfere with shots. The backboard was introduced to prevent this interference; it had the additional effect of allowing rebound shots.[11] Naismith's handwritten diaries, discovered by his granddaughter in early 2006, indicate that he was nervous about the new game he had invented, which incorporated rules from a children's game called duck on a rock, as many had failed before it.
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+ Frank Mahan, one of the players from the original first game, approached Naismith after the Christmas break, in early 1892, asking him what he intended to call his new game. Naismith replied that he hadn't thought of it because he had been focused on just getting the game started. Mahan suggested that it be called "Naismith ball", at which he laughed, saying that a name like that would kill any game. Mahan then said, "Why not call it basketball?" Naismith replied, "We have a basket and a ball, and it seems to me that would be a good name for it."[12][13] The first official game was played in the YMCA gymnasium in Albany, New York, on January 20, 1892, with nine players. The game ended at 1–0; the shot was made from 25 feet (7.6 m), on a court just half the size of a present-day Streetball or National Basketball Association (NBA) court.
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+ At the time, football was being played with 10 to a team (which was increased to 11). When winter weather got too icy to play football, teams were taken indoors, and it was convenient to have them split in half and play basketball with five on each side. By 1897–1898 teams of five became standard.
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+ Basketball's early adherents were dispatched to YMCAs throughout the United States, and it quickly spread through the United States and Canada. By 1895, it was well established at several women's high schools. While YMCA was responsible for initially developing and spreading the game, within a decade it discouraged the new sport, as rough play and rowdy crowds began to detract from YMCA's primary mission. However, other amateur sports clubs, colleges, and professional clubs quickly filled the void. In the years before World War I, the Amateur Athletic Union and the Intercollegiate Athletic Association of the United States (forerunner of the NCAA) vied for control over the rules for the game. The first pro league, the National Basketball League, was formed in 1898 to protect players from exploitation and to promote a less rough game. This league only lasted five years.
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+ James Naismith was instrumental in establishing college basketball. His colleague C.O. Beamis fielded the first college basketball team just a year after the Springfield YMCA game at the suburban Pittsburgh Geneva College.[14] Naismith himself later coached at the University of Kansas for six years, before handing the reins to renowned coach Forrest "Phog" Allen. Naismith's disciple Amos Alonzo Stagg brought basketball to the University of Chicago, while Adolph Rupp, a student of Naismith's at Kansas, enjoyed great success as coach at the University of Kentucky. On February 9, 1895, the first intercollegiate 5-on-5 game was played at Hamline University between Hamline and the School of Agriculture, which was affiliated with the University of Minnesota.[15][16][17] The School of Agriculture won in a 9–3 game.
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+ In 1901, colleges, including the University of Chicago, Columbia University, Cornell University, Dartmouth College, the University of Minnesota, the U.S. Naval Academy, the University of Colorado and Yale University began sponsoring men's games. In 1905, frequent injuries on the football field prompted President Theodore Roosevelt to suggest that colleges form a governing body, resulting in the creation of the Intercollegiate Athletic Association of the United States (IAAUS). In 1910, that body would change its name to the National Collegiate Athletic Association (NCAA). The first Canadian interuniversity basketball game was played at YMCA in Kingston, Ontario on February 6, 1904, when McGill University—Naismith's alma mater—visited Queen's University. McGill won 9–7 in overtime; the score was 7–7 at the end of regulation play, and a ten-minute overtime period settled the outcome. A good turnout of spectators watched the game.[18]
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+ The first men's national championship tournament, the National Association of Intercollegiate Basketball tournament, which still exists as the National Association of Intercollegiate Athletics (NAIA) tournament, was organized in 1937. The first national championship for NCAA teams, the National Invitation Tournament (NIT) in New York, was organized in 1938; the NCAA national tournament would begin one year later. College basketball was rocked by gambling scandals from 1948 to 1951, when dozens of players from top teams were implicated in match fixing and point shaving. Partially spurred by an association with cheating, the NIT lost support to the NCAA tournament.
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+ Before widespread school district consolidation, most American high schools were far smaller than their present-day counterparts. During the first decades of the 20th century, basketball quickly became the ideal interscholastic sport due to its modest equipment and personnel requirements. In the days before widespread television coverage of professional and college sports, the popularity of high school basketball was unrivaled in many parts of America. Perhaps the most legendary of high school teams was Indiana's Franklin Wonder Five, which took the nation by storm during the 1920s, dominating Indiana basketball and earning national recognition.
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+ Today virtually every high school in the United States fields a basketball team in varsity competition.[19] Basketball's popularity remains high, both in rural areas where they carry the identification of the entire community, as well as at some larger schools known for their basketball teams where many players go on to participate at higher levels of competition after graduation. In the 2016–17 season, 980,673 boys and girls represented their schools in interscholastic basketball competition, according to the National Federation of State High School Associations.[20] The states of Illinois, Indiana and Kentucky are particularly well known for their residents' devotion to high school basketball, commonly called Hoosier Hysteria in Indiana; the critically acclaimed film Hoosiers shows high school basketball's depth of meaning to these communities.
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+ There is currently no tournament to determine a national high school champion. The most serious effort was the National Interscholastic Basketball Tournament at the University of Chicago from 1917 to 1930. The event was organized by Amos Alonzo Stagg and sent invitations to state champion teams. The tournament started out as a mostly Midwest affair but grew. In 1929 it had 29 state champions. Faced with opposition from the National Federation of State High School Associations and North Central Association of Colleges and Schools that bore a threat of the schools losing their accreditation the last tournament was in 1930. The organizations said they were concerned that the tournament was being used to recruit professional players from the prep ranks.[21] The tournament did not invite minority schools or private/parochial schools.
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+ The National Catholic Interscholastic Basketball Tournament ran from 1924 to 1941 at Loyola University.[22] The National Catholic Invitational Basketball Tournament from 1954 to 1978 played at a series of venues, including Catholic University, Georgetown and George Mason.[23] The National Interscholastic Basketball Tournament for Black High Schools was held from 1929 to 1942 at Hampton Institute.[24] The National Invitational Interscholastic Basketball Tournament was held from 1941 to 1967 starting out at Tuskegee Institute. Following a pause during World War II it resumed at Tennessee State College in Nashville. The basis for the champion dwindled after 1954 when Brown v. Board of Education began an integration of schools. The last tournaments were held at Alabama State College from 1964 to 1967.[25]
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+ Teams abounded throughout the 1920s. There were hundreds of men's professional basketball teams in towns and cities all over the United States, and little organization of the professional game. Players jumped from team to team and teams played in armories and smoky dance halls. Leagues came and went. Barnstorming squads such as the Original Celtics and two all-African American teams, the New York Renaissance Five ("Rens") and the (still existing) Harlem Globetrotters played up to two hundred games a year on their national tours.
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+ In 1946, the Basketball Association of America (BAA) was formed. The first game was played in Toronto, Ontario, Canada between the Toronto Huskies and New York Knickerbockers on November 1, 1946. Three seasons later, in 1949, the BAA merged with the National Basketball League (NBL) to form the National Basketball Association (NBA). By the 1950s, basketball had become a major college sport, thus paving the way for a growth of interest in professional basketball. In 1959, a basketball hall of fame was founded in Springfield, Massachusetts, site of the first game. Its rosters include the names of great players, coaches, referees and people who have contributed significantly to the development of the game. The hall of fame has people who have accomplished many goals in their career in basketball. An upstart organization, the American Basketball Association, emerged in 1967 and briefly threatened the NBA's dominance until the ABA-NBA merger in 1976. Today the NBA is the top professional basketball league in the world in terms of popularity, salaries, talent, and level of competition.
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+ The NBA has featured many famous players, including George Mikan, the first dominating "big man"; ball-handling wizard Bob Cousy and defensive genius Bill Russell of the Boston Celtics; charismatic center Wilt Chamberlain, who originally played for the barnstorming Harlem Globetrotters; all-around stars Oscar Robertson and Jerry West; more recent big men Kareem Abdul-Jabbar, Shaquille O'Neal, Hakeem Olajuwon and Karl Malone; playmakers John Stockton, Isiah Thomas and Steve Nash; crowd-pleasing forwards Julius Erving and Charles Barkley; European stars Dirk Nowitzki, Pau Gasol and Tony Parker; more recent superstars LeBron James, Allen Iverson and Kobe Bryant; and the three players who many credit with ushering the professional game to its highest level of popularity during the 1980s and 1990s: Larry Bird, Earvin "Magic" Johnson, and Michael Jordan.
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+ In 2001, the NBA formed a developmental league, the National Basketball Development League (later known as the NBA D-League and then the NBA G League after a branding deal with Gatorade). As of the 2018–19 season, the G League has 27 teams.
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+ FIBA (International Basketball Federation) was formed in 1932 by eight founding nations: Argentina, Czechoslovakia, Greece, Italy, Latvia, Portugal, Romania and Switzerland. At this time, the organization only oversaw amateur players. Its acronym, derived from the French Fédération Internationale de Basket-ball Amateur, was thus "FIBA". Men's basketball was first included at the Berlin 1936 Summer Olympics, although a demonstration tournament was held in 1904. The United States defeated Canada in the first final, played outdoors. This competition has usually been dominated by the United States, whose team has won all but three titles. The first of these came in a controversial final game in Munich in 1972 against the Soviet Union, in which the ending of the game was replayed three times until the Soviet Union finally came out on top.[26] In 1950 the first FIBA World Championship for men, now known as the FIBA Basketball World Cup, was held in Argentina. Three years later, the first FIBA World Championship for women, now known as the FIBA Women's Basketball World Cup, was held in Chile. Women's basketball was added to the Olympics in 1976, which were held in Montreal, Quebec, Canada with teams such as the Soviet Union, Brazil and Australia rivaling the American squads.
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+ In 1989, FIBA allowed professional NBA players to participate in the Olympics for the first time. Prior to the 1992 Summer Olympics, only European and South American teams were allowed to field professionals in the Olympics. The United States' dominance continued with the introduction of the original Dream Team. In the 2004 Athens Olympics, the United States suffered its first Olympic loss while using professional players, falling to Puerto Rico (in a 19-point loss) and Lithuania in group games, and being eliminated in the semifinals by Argentina. It eventually won the bronze medal defeating Lithuania, finishing behind Argentina and Italy. The Redeem Team, won gold at the 2008 Olympics, and the B-Team, won gold at the 2010 FIBA World Championship in Turkey despite featuring no players from the 2008 squad. The United States continued its dominance as they won gold at the 2012 Olympics, 2014 FIBA World Cup and the 2016 Olympics.
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+ Worldwide, basketball tournaments are held for boys and girls of all age levels. The global popularity of the sport is reflected in the nationalities represented in the NBA. Players from all six inhabited continents currently play in the NBA. Top international players began coming into the NBA in the mid-1990s, including Croatians Dražen Petrović and Toni Kukoč, Serbian Vlade Divac, Lithuanians Arvydas Sabonis and Šarūnas Marčiulionis, Dutchman Rik Smits and German Detlef Schrempf.
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+ In the Philippines, the Philippine Basketball Association's first game was played on April 9, 1975 at the Araneta Coliseum in Cubao, Quezon City. Philippines. It was founded as a "rebellion" of several teams from the now-defunct Manila Industrial and Commercial Athletic Association, which was tightly controlled by the Basketball Association of the Philippines (now defunct), the then-FIBA recognized national association. Nine teams from the MICAA participated in the league's first season that opened on April 9, 1975. The NBL is Australia's pre-eminent men's professional basketball league. The league commenced in 1979, playing a winter season (April–September) and did so until the completion of the 20th season in 1998. The 1998–99 season, which commenced only months later, was the first season after the shift to the current summer season format (October–April). This shift was an attempt to avoid competing directly against Australia's various football codes. It features 8 teams from around Australia and one in New Zealand. A few players including Luc Longley, Andrew Gaze, Shane Heal, Chris Anstey and Andrew Bogut made it big internationally, becoming poster figures for the sport in Australia. The Women's National Basketball League began in 1981.
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+ Women's basketball began in 1892 at Smith College when Senda Berenson, a physical education teacher, modified Naismith's rules for women. Shortly after she was hired at Smith, she went to Naismith to learn more about the game.[27] Fascinated by the new sport and the values it could teach, she organized the first women's collegiate basketball game on March 21, 1893, when her Smith freshmen and sophomores played against one another.[28] However, the first women's interinstitutional game was played in 1892 between the University of California and Miss Head's School.[29] Berenson's rules were first published in 1899, and two years later she became the editor of A. G. Spalding's first Women's Basketball Guide.[28] Berenson's freshmen played the sophomore class in the first women's intercollegiate basketball game at Smith College, March 21, 1893.[30] The same year, Mount Holyoke and Sophie Newcomb College (coached by Clara Gregory Baer) women began playing basketball. By 1895, the game had spread to colleges across the country, including Wellesley, Vassar, and Bryn Mawr. The first intercollegiate women's game was on April 4, 1896. Stanford women played Berkeley, 9-on-9, ending in a 2–1 Stanford victory.
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+ Women's basketball development was more structured than that for men in the early years. In 1905, the Executive Committee on Basket Ball Rules (National Women's Basketball Committee) was created by the American Physical Education Association.[31] These rules called for six to nine players per team and 11 officials. The International Women's Sports Federation (1924) included a women's basketball competition. 37 women's high school varsity basketball or state tournaments were held by 1925. And in 1926, the Amateur Athletic Union backed the first national women's basketball championship, complete with men's rules.[31] The Edmonton Grads, a touring Canadian women's team based in Edmonton, Alberta, operated between 1915 and 1940. The Grads toured all over North America, and were exceptionally successful. They posted a record of 522 wins and only 20 losses over that span, as they met any team that wanted to challenge them, funding their tours from gate receipts.[32] The Grads also shone on several exhibition trips to Europe, and won four consecutive exhibition Olympics tournaments, in 1924, 1928, 1932, and 1936; however, women's basketball was not an official Olympic sport until 1976. The Grads' players were unpaid, and had to remain single. The Grads' style focused on team play, without overly emphasizing skills of individual players. The first women's AAU All-America team was chosen in 1929.[31] Women's industrial leagues sprang up throughout the United States, producing famous athletes, including Babe Didrikson of the Golden Cyclones, and the All American Red Heads Team, which competed against men's teams, using men's rules. By 1938, the women's national championship changed from a three-court game to two-court game with six players per team.[31]
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+ The NBA-backed Women's National Basketball Association (WNBA) began in 1997. Though it had shaky attendance figures, several marquee players (Lisa Leslie, Diana Taurasi, and Candace Parker among others) have helped the league's popularity and level of competition. Other professional women's basketball leagues in the United States, such as the American Basketball League (1996–98), have folded in part because of the popularity of the WNBA. The WNBA has been looked at by many as a niche league. However, the league has recently taken steps forward. In June 2007, the WNBA signed a contract extension with ESPN. The new television deal ran from 2009 to 2016. Along with this deal, came the first ever rights fees to be paid to a women's professional sports league. Over the eight years of the contract, "millions and millions of dollars" were "dispersed to the league's teams." In a March 12, 2009 article, NBA commissioner David Stern said that in the bad economy, "the NBA is far less profitable than the WNBA. We're losing a lot of money among a large number of teams. We're budgeting the WNBA to break even this year."[33]
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+ Measurements and time limits discussed in this section often vary among tournaments and organizations; international and NBA rules are used in this section.
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+ The object of the game is to outscore one's opponents by throwing the ball through the opponents' basket from above while preventing the opponents from doing so on their own. An attempt to score in this way is called a shot. A successful shot is worth two points, or three points if it is taken from beyond the three-point arc 6.75 metres (22 ft 2 in) from the basket in international games[citation needed] and 23 feet 9 inches (7.24 m) in NBA games.[34] A one-point shot can be earned when shooting from the foul line after a foul is made. After a team has scored from a field goal or free throw, play is resumed with a throw-in awarded to the non-scoring team taken from a point beyond the endline of the court where the points(s) were scored.[35]
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+ Games are played in four quarters of 10 (FIBA)[36] or 12 minutes (NBA).[37] College men's games use two 20-minute halves,[38] college women's games use 10-minute quarters,[39] and most United States high school varsity games use 8-minute quarters; however, this varies from state to state.[40][41] 15 minutes are allowed for a half-time break under FIBA, NBA, and NCAA rules[38][42][43] and 10 minutes in United States high schools.[40] Overtime periods are five minutes in length[38][44][45] except for high school, which is four minutes in length.[40] Teams exchange baskets for the second half. The time allowed is actual playing time; the clock is stopped while the play is not active. Therefore, games generally take much longer to complete than the allotted game time, typically about two hours.
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+ Five players from each team may be on the court at one time.[46][47][48][49] Substitutions are unlimited but can only be done when play is stopped. Teams also have a coach, who oversees the development and strategies of the team, and other team personnel such as assistant coaches, managers, statisticians, doctors and trainers.
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+ For both men's and women's teams, a standard uniform consists of a pair of shorts and a jersey with a clearly visible number, unique within the team, printed on both the front and back. Players wear high-top sneakers that provide extra ankle support. Typically, team names, players' names and, outside of North America, sponsors are printed on the uniforms.
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+ A limited number of time-outs, clock stoppages requested by a coach (or sometimes mandated in the NBA) for a short meeting with the players, are allowed. They generally last no longer than one minute (100 seconds in the NBA) unless, for televised games, a commercial break is needed.
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+ The game is controlled by the officials consisting of the referee (referred to as crew chief in the NBA), one or two umpires (referred to as referees in the NBA) and the table officials. For college, the NBA, and many high schools, there are a total of three referees on the court. The table officials are responsible for keeping track of each team's scoring, timekeeping, individual and team fouls, player substitutions, team possession arrow, and the shot clock.
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+ The only essential equipment in a basketball game is the ball and the court: a flat, rectangular surface with baskets at opposite ends. Competitive levels require the use of more equipment such as clocks, score sheets, scoreboard(s), alternating possession arrows, and whistle-operated stop-clock systems.
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+ A regulation basketball court in international games is 91.9 feet (28.0 meters) long and 49.2 feet (15 meters) wide. In the NBA and NCAA the court is 94 by 50 feet (29 by 15 meters).[34] Most courts have wood flooring, usually constructed from maple planks running in the same direction as the longer court dimension.[50][51] The name and logo of the home team is usually painted on or around the center circle.
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+ The basket is a steel rim 18 inches (46 cm) diameter with an attached net affixed to a backboard that measures 6 by 3.5 feet (1.8 by 1.1 meters) and one basket is at each end of the court. The white outlined box on the backboard is 18 inches (46 cm) high and 2 feet (61 cm) wide. At almost all levels of competition, the top of the rim is exactly 10 feet (3.05 meters) above the court and 4 feet (1.22 meters) inside the baseline. While variation is possible in the dimensions of the court and backboard, it is considered important for the basket to be of the correct height – a rim that is off by just a few inches can have an adverse effect on shooting.
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+
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+ The size of the basketball is also regulated. For men, the official ball is 29.5 inches (75 cm) in circumference (size 7, or a "295 ball") and weighs 22 oz (623.69 grams). If women are playing, the official basketball size is 28.5 inches (72 cm) in circumference (size 6, or a "285 ball") with a weight of 20 oz (567 grams). In 3x3, a formalized version of the halfcourt 3-on-3 game, a dedicated ball with the circumference of a size 6 ball but the weight of a size 7 ball is used in all competitions (men's, women's, and mixed teams).[52]
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+
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+ The ball may be advanced toward the basket by being shot, passed between players, thrown, tapped, rolled or dribbled (bouncing the ball while running).
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+ The ball must stay within the court; the last team to touch the ball before it travels out of bounds forfeits possession. The ball is out of bounds if it touches a boundary line, or touches any player or object that is out of bounds.
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+
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+ There are limits placed on the steps a player may take without dribbling, which commonly results in an infraction known as traveling. Nor may a player stop his dribble and then resume dribbling. A dribble that touches both hands is considered stopping the dribble, giving this infraction the name double dribble. Within a dribble, the player cannot carry the ball by placing his hand on the bottom of the ball; doing so is known as carrying the ball. A team, once having established ball control in the front half of their court, may not return the ball to the backcourt and be the first to touch it. A violation of these rules results in loss of possession.
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+ The ball may not be kicked, nor be struck with the fist. For the offense, a violation of these rules results in loss of possession; for the defense, most leagues reset the shot clock and the offensive team is given possession of the ball out of bounds.
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+
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+ There are limits imposed on the time taken before progressing the ball past halfway (8 seconds in FIBA and the NBA; 10 seconds in NCAA and high school for both sexes), before attempting a shot (24 seconds in FIBA, the NBA, and U Sports (Canadian universities) play for both sexes, and 30 seconds in NCAA play for both sexes), holding the ball while closely guarded (5 seconds), and remaining in the restricted area known as the free-throw lane, (or the "key") (3 seconds). These rules are designed to promote more offense.
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+
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+ Basket interference, or goaltending is a violation charged when a player illegally interferes with a shot. This violation is incurred when a player touches the ball on its downward trajectory to the basket, unless it is obvious that the ball has no chance of entering the basket, if a player touches the ball while it is in the rim, or in the area extended upwards from the basket, or if a player reaches through the basket to interfere with the shot. When a defensive player is charged with goaltending, the basket is awarded. If an offensive player commits the infraction, the basket is cancelled. In either case possession of the ball is turned over to the defensive team.
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+
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+ An attempt to unfairly disadvantage an opponent through certain types of physical contact is illegal and is called a personal foul. These are most commonly committed by defensive players; however, they can be committed by offensive players as well. Players who are fouled either receive the ball to pass inbounds again, or receive one or more free throws if they are fouled in the act of shooting, depending on whether the shot was successful. One point is awarded for making a free throw, which is attempted from a line 15 feet (4.6 m) from the basket.
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+ The referee is responsible for judging whether contact is illegal, sometimes resulting in controversy. The calling of fouls can vary between games, leagues and referees.
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+ There is a second category of fouls called technical fouls, which may be charged for various rules violations including failure to properly record a player in the scorebook, or for unsportsmanlike conduct. These infractions result in one or two free throws, which may be taken by any of the five players on the court at the time. Repeated incidents can result in disqualification. A blatant foul involving physical contact that is either excessive or unnecessary is called an intentional foul (flagrant foul in the NBA). In FIBA and NCAA women's basketball, a foul resulting in ejection is called a disqualifying foul, while in leagues other than the NBA, such a foul is referred to as flagrant.
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+
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+ If a team exceeds a certain limit of team fouls in a given period (quarter or half) – four for NBA, NCAA women's, and international games – the opposing team is awarded one or two free throws on all subsequent non-shooting fouls for that period, the number depending on the league. In the US college men's game and high school games for both sexes, if a team reaches 7 fouls in a half, the opposing team is awarded one free throw, along with a second shot if the first is made. This is called shooting "one-and-one". If a team exceeds 10 fouls in the half, the opposing team is awarded two free throws on all subsequent fouls for the half.
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+ When a team shoots foul shots, the opponents may not interfere with the shooter, nor may they try to regain possession until the last or potentially last free throw is in the air.
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+ After a team has committed a specified number of fouls, the other team is said to be "in the bonus". On scoreboards, this is usually signified with an indicator light reading "Bonus" or "Penalty" with an illuminated directional arrow or dot indicating that team is to receive free throws when fouled by the opposing team. (Some scoreboards also indicate the number of fouls committed.)
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+ If a team misses the first shot of a two-shot situation, the opposing team must wait for the completion of the second shot before attempting to reclaim possession of the ball and continuing play.
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+ If a player is fouled while attempting a shot and the shot is unsuccessful, the player is awarded a number of free throws equal to the value of the attempted shot. A player fouled while attempting a regular two-point shot thus receives two shots, and a player fouled while attempting a three-point shot receives three shots.
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+ If a player is fouled while attempting a shot and the shot is successful, typically the player will be awarded one additional free throw for one point. In combination with a regular shot, this is called a "three-point play" or "four-point play" (or more colloquially, an "and one") because of the basket made at the time of the foul (2 or 3 points) and the additional free throw (1 point).
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+ Although the rules do not specify any positions whatsoever, they have evolved as part of basketball. During the early years of basketball's evolution, two guards, two forwards, and one center were used. In more recent times specific positions evolved, but the current trend, advocated by many top coaches including Mike Krzyzewski is towards positionless basketball, where big guys are free to shoot from outside and dribble if their skill allows it.[53] Popular descriptions of positions include:
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+ Point guard (often called the "1") : usually the fastest player on the team, organizes the team's offense by controlling the ball and making sure that it gets to the right player at the right time.
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+ Shooting guard (the "2") : creates a high volume of shots on offense, mainly long-ranged; and guards the opponent's best perimeter player on defense.
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+ Small forward (the "3") : often primarily responsible for scoring points via cuts to the basket and dribble penetration; on defense seeks rebounds and steals, but sometimes plays more actively.
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+ Power forward (the "4"): plays offensively often with their back to the basket; on defense, plays under the basket (in a zone defense) or against the opposing power forward (in man-to-man defense).
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+ Center (the "5"): uses height and size to score (on offense), to protect the basket closely (on defense), or to rebound.
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+ The above descriptions are flexible. For most teams today, the shooting guard and small forward have very similar responsibilities and are often called the wings, as do the power forward and center, who are often called post players. While most teams describe two players as guards, two as forwards, and one as a center, on some occasions teams choose to call them by different designations.
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+ There are two main defensive strategies: zone defense and man-to-man defense. In a zone defense, each player is assigned to guard a specific area of the court. Zone defenses often allow the defense to double team the ball, a manoeuver known as a trap. In a man-to-man defense, each defensive player guards a specific opponent.
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+ Offensive plays are more varied, normally involving planned passes and movement by players without the ball. A quick movement by an offensive player without the ball to gain an advantageous position is known as a cut. A legal attempt by an offensive player to stop an opponent from guarding a teammate, by standing in the defender's way such that the teammate cuts next to him, is a screen or pick. The two plays are combined in the pick and roll, in which a player sets a pick and then "rolls" away from the pick towards the basket. Screens and cuts are very important in offensive plays; these allow the quick passes and teamwork, which can lead to a successful basket. Teams almost always have several offensive plays planned to ensure their movement is not predictable. On court, the point guard is usually responsible for indicating which play will occur.
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+ Shooting is the act of attempting to score points by throwing the ball through the basket, methods varying with players and situations.
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+ Typically, a player faces the basket with both feet facing the basket. A player will rest the ball on the fingertips of the dominant hand (the shooting arm) slightly above the head, with the other hand supporting the side of the ball. The ball is usually shot by jumping (though not always) and extending the shooting arm. The shooting arm, fully extended with the wrist fully bent, is held stationary for a moment following the release of the ball, known as a follow-through. Players often try to put a steady backspin on the ball to absorb its impact with the rim. The ideal trajectory of the shot is somewhat controversial, but generally a proper arc is recommended. Players may shoot directly into the basket or may use the backboard to redirect the ball into the basket.
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+ The two most common shots that use the above described setup are the set shot and the jump shot. Both are preceeded by a crouching action which preloads the muscles and increases the power of the shot. In a set shot the shooter straightens up and throws from a standing position with neither foot leaving the floor; this is typically used for free throws. For a jump shot, the throw is taken in mid-air with the ball being released near the top of the jump. This provides much greater power and range, and it also allows the player to elevate over the defender. Failure to release the ball before the feet return to the floor is considered a traveling violation.
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+ Another common shot is called the lay-up. This shot requires the player to be in motion toward the basket, and to "lay" the ball "up" and into the basket, typically off the backboard (the backboard-free, underhand version is called a finger roll). The most crowd-pleasing and typically highest-percentage accuracy shot is the slam dunk, in which the player jumps very high and throws the ball downward, through the basket while touching it.
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+ Another shot that is becoming common[citation needed] is the "circus shot". The circus shot is a low-percentage shot that is flipped, heaved, scooped, or flung toward the hoop while the shooter is off-balance, airborne, falling down, and/or facing away from the basket. A back-shot is a shot taken when the player is facing away from the basket, and may be shot with the dominant hand, or both; but there is a very low chance that the shot will be successful.
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+ A shot that misses both the rim and the backboard completely is referred to as an air ball. A particularly bad shot, or one that only hits the backboard, is jocularly called a brick. The hang time is the length of time a player stays in the air after jumping, either to make a slam dunk, lay-up or jump shot.
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+ The objective of rebounding is to successfully gain possession of the basketball after a missed field goal or free throw, as it rebounds from the hoop or backboard. This plays a major role in the game, as most possessions end when a team misses a shot. There are two categories of rebounds: offensive rebounds, in which the ball is recovered by the offensive side and does not change possession, and defensive rebounds, in which the defending team gains possession of the loose ball. The majority of rebounds are defensive, as the team on defense tends to be in better position to recover missed shots.
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+ A pass is a method of moving the ball between players. Most passes are accompanied by a step forward to increase power and are followed through with the hands to ensure accuracy.
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+ A staple pass is the chest pass. The ball is passed directly from the passer's chest to the receiver's chest. A proper chest pass involves an outward snap of the thumbs to add velocity and leaves the defence little time to react.
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+ Another type of pass is the bounce pass. Here, the passer bounces the ball crisply about two-thirds of the way from his own chest to the receiver. The ball strikes the court and bounces up toward the receiver. The bounce pass takes longer to complete than the chest pass, but it is also harder for the opposing team to intercept (kicking the ball deliberately is a violation). Thus, players often use the bounce pass in crowded moments, or to pass around a defender.
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+ The overhead pass is used to pass the ball over a defender. The ball is released while over the passer's head.
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+ The outlet pass occurs after a team gets a defensive rebound. The next pass after the rebound is the outlet pass.
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+ The crucial aspect of any good pass is it being difficult to intercept. Good passers can pass the ball with great accuracy and they know exactly where each of their other teammates prefers to receive the ball. A special way of doing this is passing the ball without looking at the receiving teammate. This is called a no-look pass.
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+ Another advanced style of passing is the behind-the-back pass, which, as the description implies, involves throwing the ball behind the passer's back to a teammate. Although some players can perform such a pass effectively, many coaches discourage no-look or behind-the-back passes, believing them to be difficult to control and more likely to result in turnovers or violations.
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+ Dribbling is the act of bouncing the ball continuously with one hand and is a requirement for a player to take steps with the ball. To dribble, a player pushes the ball down towards the ground with the fingertips rather than patting it; this ensures greater control.
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+ When dribbling past an opponent, the dribbler should dribble with the hand farthest from the opponent, making it more difficult for the defensive player to get to the ball. It is therefore important for a player to be able to dribble competently with both hands.
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+ Good dribblers (or "ball handlers") tend to bounce the ball low to the ground, reducing the distance of travel of the ball from the floor to the hand, making it more difficult for the defender to "steal" the ball. Good ball handlers frequently dribble behind their backs, between their legs, and switch directions suddenly, making a less predictable dribbling pattern that is more difficult to defend against. This is called a crossover, which is the most effective way to move past defenders while dribbling.
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+ A skilled player can dribble without watching the ball, using the dribbling motion or peripheral vision to keep track of the ball's location. By not having to focus on the ball, a player can look for teammates or scoring opportunities, as well as avoid the danger of having someone steal the ball away from him/her.
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+ A block is performed when, after a shot is attempted, a defender succeeds in altering the shot by touching the ball. In almost all variants of play, it is illegal to touch the ball after it is in the downward path of its arc; this is known as goaltending. It is also illegal under NBA and Men's NCAA basketball to block a shot after it has touched the backboard, or when any part of the ball is directly above the rim. Under international rules it is illegal to block a shot that is in the downward path of its arc or one that has touched the backboard until the ball has hit the rim. After the ball hits the rim, it is again legal to touch it even though it is no longer considered as a block performed.
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+ To block a shot, a player has to be able to reach a point higher than where the shot is released. Thus, height can be an advantage in blocking. Players who are taller and playing the power forward or center positions generally record more blocks than players who are shorter and playing the guard positions. However, with good timing and a sufficiently high vertical leap, even shorter players can be effective shot blockers.
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+ At the professional level, most male players are above 6 feet 3 inches (1.91 m) and most women above 5 feet 7 inches (1.70 m). Guards, for whom physical coordination and ball-handling skills are crucial, tend to be the smallest players. Almost all forwards in the top men's pro leagues are 6 feet 6 inches (1.98 m) or taller. Most centers are over 6 feet 10 inches (2.08 m) tall. According to a survey given to all NBA teams,[when?] the average height of all NBA players is just under 6 feet 7 inches (2.01 m), with the average weight being close to 222 pounds (101 kg). The tallest players ever in the NBA were Manute Bol and Gheorghe Mureșan, who were both 7 feet 7 inches (2.31 m) tall. At 7 feet 2 inches (2.18 m), Margo Dydek was the tallest player in the history of the WNBA.
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+ The shortest player ever to play in the NBA is Muggsy Bogues at 5 feet 3 inches (1.60 m).[54] Other short players have thrived at the pro level. Anthony "Spud" Webb was just 5 feet 7 inches (1.70 m) tall, but had a 42-inch (1.1 m) vertical leap, giving him significant height when jumping. While shorter players are often at a disadvantage in certain aspects of the game, their ability to navigate quickly through crowded areas of the court and steal the ball by reaching low are strengths.
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+ Players regularly inflate their height. Many prospects exaggerate their height while in high school or college to make themselves more appealing to coaches and scouts, who prefer taller players. Charles Barkley stated; "I've been measured at 6-5, 6-4 ​3⁄4. But I started in college at 6-6." Sam Smith, a former writer from The Chicago Tribune, said: "We sort of know the heights, because after camp, the sheet comes out. But you use that height, and the player gets mad. And then you hear from his agent. Or you file your story with the right height, and the copy desk changes it because they have the 'official' N.B.A. media guide, which is wrong. So you sort of go along with the joke."[55] In the NBA, there is no standard on whether a player's listed height uses their measurement with shoes on or without. The NBA Draft Combine, which most players attend before the draft, provides both measurements. Thereafter, a player's team is solely responsible for their listed height, which can vary depending on the process selected.[56][57]
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+ Notable players who overstated their height include:
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+ On rare occasions, some players will understate their actual heights, not to be repositioned. One example is Kevin Durant, whose listed height is 6 feet 9 inches (2.06 m), while his actual height is 7 feet 0 inches (2.13 m). Durant's reasoning was, "Really, that's the prototypical size for a small forward. Anything taller than that, and they'll start saying, 'Ah, he's a power forward."[62]
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+ Variations of basketball are activities based on the game of basketball, using common basketball skills and equipment (primarily the ball and basket). Some variations are only superficial rules changes, while others are distinct games with varying degrees of basketball influences. Other variations include children's games, contests or activities meant to help players reinforce skills.
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+ There are principal basketball sports with variations on basketball including Wheelchair basketball, Water basketball, Beach basketball, Slamball, Streetball and Unicycle basketball. An earlier version of basketball, played primarily by women and girls, was Six-on-six basketball. Horseball is a game played on horseback where a ball is handled and points are scored by shooting it through a high net (approximately 1.5m×1.5m). The sport is like a combination of polo, rugby, and basketball. There is even a form played on donkeys known as Donkey basketball, but that version has come under attack from animal rights groups.
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+ Perhaps the single most common variation of basketball is the half-court game, played in informal settings without referees or strict rules. Only one basket is used, and the ball must be "taken back" or "cleared" – passed or dribbled outside the three-point line each time possession of the ball changes from one team to the other. Half-court games require less cardiovascular stamina, since players need not run back and forth a full court. Half-court raises the number of players that can use a court or, conversely, can be played if there is an insufficient number to form full 5-on-5 teams.
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+ Half-court basketball is usually played 1-on-1, 2-on-2 or 3-on-3. The latter variation is gradually gaining official recognition as 3x3, originally known as FIBA 33. It was first tested at the 2007 Asian Indoor Games in Macau and the first official tournaments were held at the 2009 Asian Youth Games and the 2010 Youth Olympics, both in Singapore. The first FIBA 3x3 Youth World Championships[63] were held in Rimini, Italy in 2011, with the first FIBA 3x3 World Championships for senior teams following a year later in Athens. The sport is highly tipped to become an Olympic sport as early as 2016.[64] In the summer of 2017, the BIG3 basketball league, a professional 3x3 half court basketball league that features former NBA players, began. The BIG3 features several rule variants including a four-point field goal.[65]
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+ There are also other basketball sports, such as:
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+ Spin-offs from basketball that are now separate sports include:
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+ Basketball has been adopted by various social groups, which have established their own environments and sometimes their own rules. Such socialized forms of basketball include the following.
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+ Basketball is played widely casually in schools and colleges where fun, entertainment and camaraderie rule rather than winning a game.
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+ Disabled basketball is played by various disabled groups, such as the deaf and physically crippled people.
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+ Show basketball is performed by entertainment basketball show teams, the prime example being the Harlem Globetrotters. There are even specialized entertainment teams, such as teams of celebrities, people with short heights and others.
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+ Fantasy basketball was popularized during the 1990s after the advent of the Internet. Those who play this game are sometimes referred to as General Managers, who draft actual NBA players and compute their basketball statistics. The game was popularized by ESPN Fantasy Sports, NBA.com, and Yahoo! Fantasy Sports. Other sports websites provided the same format keeping the game interesting with participants actually owning specific players.
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1
+ The Twelve Labours of Heracles or Hercules (Greek: οἱ Ἡρακλέους ἆθλοι, hoi Hērakléous âthloi)[1][2] are a series of episodes concerning a penance carried out by Heracles, the greatest of the Greek heroes, whose name was later romanised as Hercules. They were accomplished at the service of King Eurystheus. The episodes were later connected by a continuous narrative. The establishment of a fixed cycle of twelve labours was attributed by the Greeks to an epic poem, now lost, written by Peisander, dated about 600 BC.[3] After Heracles killed his wife and children, he went to the oracle at Delphi. He prayed to the god Apollo for guidance. Heracles was told to serve the king of Mycenae, Eurystheus, for ten years. During this time, he is sent to perform a series of difficult feats, called labours.[4]
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+ Driven mad by Hera (queen of the gods), Heracles slew his sons by his wife Megara.[5] After recovering his sanity, Heracles deeply regretted his actions; he was purified by King Thespius, then traveled to Delphi to inquire how he could atone for his actions. Pythia, the Oracle of Delphi, advised him to go to Tiryns and serve his cousin, King Eurystheus, for ten years, performing whatever labours Eurystheus might set him; in return, he would be rewarded with immortality. Heracles despaired at this, loathing to serve a man whom he knew to be far inferior to himself, yet fearing to oppose his father Zeus. Eventually, he placed himself at Eurystheus's disposal.
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+ Eurystheus originally ordered Heracles to perform ten labours. Heracles accomplished these tasks, but Eurystheus refused to recognize two: the slaying of the Lernaean Hydra, as Heracles' nephew and charioteer Iolaus had helped him; and the cleansing of the Augeas, because Heracles accepted payment for the labour. Eurystheus set two more tasks (fetching the Golden Apples of Hesperides and capturing Cerberus), which Heracles also performed, bringing the total number of tasks to twelve.
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+ As they survive, the labours of Heracles are not recounted in any single place, but must be reassembled from many sources. Ruck and Staples[6] assert that there is no one way to interpret the labours, but that six were located in the Peloponnese, culminating with the rededication of Olympia. Six others took the hero farther afield, to places that were, per Ruck, "all previously strongholds of Hera or the 'Goddess' and were Entrances to the Netherworld".[6] In each case, the pattern was the same: Heracles was sent to kill or subdue, or to fetch back for Eurystheus (as Hera's representative) a magical animal or plant.
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+ A famous depiction of the labours in Greek sculpture is found on the metopes of the Temple of Zeus at Olympia, which date to the 450s BC.[citation needed]
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+ In his labours, Heracles was sometimes accompanied by a male companion (an eromenos), according to Licymnius[citation needed] and others, such as Iolaus, his nephew. Although he was supposed to perform only ten labours, this assistance led to two labours being disqualified: Eurystheus refused to recognize slaying the Hydra, because Iolaus helped him, and the cleansing of the Augean stables, because Heracles was paid for his services and/or because the rivers did the work. Several of the labours involved the offspring (by various accounts) of Typhon and his mate Echidna, all overcome by Heracles.
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+ A traditional order of the labours found in the Bibliotheca[7] by Pseudo-Apollodorus is:
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+ Heracles wandered the area until he came to the town of Cleonae. There he met a boy who said that if Heracles slew the Nemean lion and returned alive within 30 days, the town would sacrifice a lion to Zeus, but if he did not return within 30 days or if he died, the boy would sacrifice himself to Zeus. Another version claims that he met Molorchos, a shepherd who had lost his son to the lion, saying that if he came back within 30 days, a ram would be sacrificed to Zeus. If he did not return within 30 days, it would be sacrificed to the dead Heracles as a mourning offering.
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+ While searching for the lion, Heracles fletched some arrows to use against it, not knowing that its golden fur was impenetrable. When he found and shot the lion, firing at it with his bow, he discovered the fur's protective property as the arrow bounced harmlessly off the creature's thigh. After some time, Heracles made the lion return to his cave. The cave had two entrances, one of which Heracles blocked; he then entered the other. In those dark and close quarters, Heracles stunned the beast with his club and, using his immense strength, strangled it to death. During the fight the lion bit off one of his fingers.[8] Others say that he shot arrows at it, eventually shooting it in the unarmored mouth. After slaying the lion, he tried to skin it with a knife from his belt, but failed. He then tried sharpening the knife with a stone and even tried with the stone itself. Finally, Athena, noticing the hero's plight, told Heracles to use one of the lion's own claws to skin the pelt. Others say that Heracles' armor was, in fact, the hide of the Lion of Cithaeron.
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+ When he returned on the 30th day carrying the carcass of the lion on his shoulders, King Eurystheus was amazed and terrified. Eurystheus forbade him ever again to enter the city; from then on he was to display the fruits of his labours outside the city gates. Eurystheus would then tell Heracles his tasks through a herald, not personally. Eurystheus even had a large bronze jar made for him in which to hide from Heracles if need be. Eurystheus then warned him that the tasks would become increasingly difficult.
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+ Heracles' second labour was to slay the Lernaean Hydra, which Hera had raised just to slay Heracles. Upon reaching the swamp near Lake Lerna, where the Hydra dwelt, Heracles used a cloth to cover his mouth and nose to protect himself from the poisonous fumes. He fired flaming arrows into the Hydra's lair, the spring of Amymone, a deep cave that it only came out of to terrorize neighboring villages.[9] He then confronted the Hydra, wielding a harvesting sickle (according to some early vase-paintings), a sword or his famed club. Ruck and Staples (1994: 170) have pointed out that the chthonic creature's reaction was botanical: upon cutting off each of its heads he found that two grew back, an expression of the hopelessness of such a struggle for any but the hero. Additionally, one of the Hydra's heads - the middle one - was immortal.
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+ The details of the struggle are explicit in the Bibliotheca (2.5.2): realizing that he could not defeat the Hydra in this way, Heracles called on his nephew Iolaus for help. His nephew then came upon the idea (possibly inspired by Athena) of using a firebrand to scorch the neck stumps after each decapitation. Heracles cut off each head and Iolaus cauterized the open stumps. Seeing that Heracles was winning the struggle, Hera sent a giant crab to distract him. He crushed it under his mighty foot. He cut off the Hydra's one immortal head with a golden sword given to him by Athena. Heracles placed it under a great rock on the sacred way between Lerna and Elaius (Kerenyi 1959:144), and dipped his arrows in the Hydra's poisonous blood, and so his second task was complete. The alternative version of this myth is that after cutting off one head, he then dipped his sword in it and used its venom to burn each head so it could not grow back. Hera, upset that Heracles had slain the beast she raised to kill him, placed it in the dark blue vault of the sky as the constellation Hydra. She then turned the crab into the constellation Cancer.
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+ Later, Heracles used an arrow dipped in the Hydra's poisonous blood to kill the centaur Nessus; and Nessus's tainted blood was applied to the Tunic of Nessus, by which the centaur had his posthumous revenge. Both Strabo and Pausanias report that the stench of the river Anigrus in Elis, making all the fish of the river inedible, was reputed to be due to the Hydra's venom, washed from the arrows Heracles used on the centaur.[10]
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+ Eurystheus and Hera were greatly angered that Heracles had survived the Nemean Lion and the Lernaean Hydra. For the third labour, they found a task which they thought would spell doom for the hero. It was not slaying a beast or monster, as it had already been established that Heracles could overcome even the most fearsome opponents. Instead, Eurystheus ordered him to capture the Ceryneian Hind, which was so fast that it could outrun an arrow.
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+ After beginning the search, Heracles awoke from sleeping and saw the hind by the glint on its antlers. Heracles then chased the hind on foot for a full year through Greece, Thrace, Istria, and the land of the Hyperboreans. In some versions, he captured the hind while it slept, rendering it lame with a trap net. In other versions, he encountered Artemis in her temple; she told him to leave the hind and tell Eurystheus all that had happened, and his third labour would be considered to be completed. Yet another version claims that Heracles trapped the Hind with an arrow between its forelegs.
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+ Eurystheus had given Heracles this task hoping to incite Artemis' anger at Heracles for his desecration of her sacred animal. As he was returning with the hind, Heracles encountered Artemis and her brother Apollo. He begged the goddess for forgiveness, explaining that he had to catch it as part of his penance, but he promised to return it. Artemis forgave him, foiling Eurystheus' plan to have her punish him.
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+ Upon bringing the hind to Eurystheus, he was told that it was to become part of the King's menagerie. Heracles knew that he had to return the hind as he had promised, so he agreed to hand it over on the condition that Eurystheus himself come out and take it from him. The King came out, but the moment that Heracles let the hind go, it sprinted back to its mistress and Heracles left, saying that Eurystheus had not been quick enough.
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+ Eurystheus was disappointed that Heracles had overcome yet another creature and was humiliated by the hind's escape, so he assigned Heracles another dangerous task. By some accounts, the fourth labour was to bring the fearsome Erymanthian Boar back to Eurystheus alive (there is no single definitive telling of the labours). On the way to Mount Erymanthos where the boar lived, Heracles visited Pholus ("caveman"), a kind and hospitable centaur and old friend. Heracles ate with Pholus in his cavern (though the centaur devoured his meat raw) and asked for wine. Pholus had only one jar of wine, a gift from Dionysus to all the centaurs on Mount Erymanthos. Heracles convinced him to open it, and the smell attracted the other centaurs. They did not understand that wine needs to be tempered with water, became drunk, and attacked Heracles. Heracles shot at them with his poisonous arrows, killing many, and the centaurs retreated all the way to Chiron's cave.
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+ Pholus was curious why the arrows caused so much death. He picked one up but dropped it, and the arrow stabbed his hoof, poisoning him. One version states that a stray arrow hit Chiron as well. He was immortal, but he still felt the pain. Chiron's pain was so great that he volunteered to give up his immortality and take the place of Prometheus, who had been chained to the top of a mountain to have his liver eaten daily by an eagle. Prometheus' torturer, the eagle, continued its torture on Chiron, so Heracles shot it dead with an arrow. It is generally accepted that the tale was meant to show Heracles as being the recipient of Chiron's surrendered immortality. However, this tale contradicts the fact that Chiron later taught Achilles. The tale of the centaurs sometimes appears in other parts of the twelve labours, as does the freeing of Prometheus.
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+ Heracles had visited Chiron to gain advice on how to catch the boar, and Chiron had told him to drive it into thick snow, which sets this labour in mid-winter. Heracles caught the boar, bound it, and carried it back to Eurystheus, who was frightened of it and ducked down in his half-buried storage pithos, begging Heracles to get rid of the beast.
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+ The fifth labour was to clean the stables of King Augeas. This assignment was intended to be both humiliating (rather than impressive, as the previous labours had been) and impossible, since the livestock were divinely healthy (and immortal) and therefore produced an enormous quantity of dung. The Augean Stables (/ɔːˈdʒiːən/) had not been cleaned in over 30 years, and over 1,000 cattle lived there. However, Heracles succeeded by re-routing the rivers Alpheus and Peneus to wash out the filth.
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+ Before starting on the task, Heracles had asked Augeas for one-tenth of the cattle if he finished the task in one day, and Augeas agreed. But afterwards Augeas refused to honour the agreement on the grounds that Heracles had been ordered to carry out the task by Eurystheus anyway. Heracles claimed his reward in court, and was supported by Augeas' son Phyleus. Augeas banished them both before the court had ruled. Heracles returned, slew Augeas, and gave his kingdom to Phyleus. Heracles then founded the Olympic Games.
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+ The success of this labour was ultimately discounted as the rushing waters had done the work of cleaning the stables and because Heracles was paid for doing the labour.
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+ Eurystheus said that Heracles still had seven labours to perform.[11]
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+
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+ The sixth labour was to defeat the Stymphalian birds, man-eating birds with beaks made of bronze and sharp metallic feathers they could launch at their victims. They were sacred to Ares, the god of war. Furthermore, their dung was highly toxic. They had migrated to Lake Stymphalia in Arcadia, where they bred quickly and took over the countryside, destroying local crops, fruit trees, and townspeople. Heracles could not go too far into the swamp, for it would not support his weight. Athena, noticing the hero's plight, gave Heracles a rattle which Hephaestus had made especially for the occasion. Heracles shook the rattle and frightened the birds into the air. Heracles then shot many of them with his arrows. The rest flew far away, never to return. The Argonauts would later encounter them.
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+ The seventh labour was to capture the Cretan Bull, father of the Minotaur. Heracles sailed to Crete, where King Minos gave Heracles permission to take the bull away and even offered him assistance (which Heracles declined, plausibly because he did not want the labour to be discounted as before).[12] The bull had been wreaking havoc on Crete by uprooting crops and leveling orchard walls. Heracles sneaked up behind the bull and then used his hands to throttle it (stopping before it was killed), and then shipped it back to Tiryns. Eurystheus, who hid in his pithos at first sight of the creature, wanted to sacrifice the bull to Hera, who hated Heracles. She refused the sacrifice because it reflected glory on Heracles. The bull was released and wandered into Marathon, becoming known as the Marathonian Bull.[12] Theseus would later sacrifice the bull to Athena and/or Apollo.
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+
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+ As the eighth of his Twelve Labours, also categorised as the second of the Non-Peloponneisan labours,[13] Heracles was sent by King Eurystheus to steal the Mares from Diomedes. The mares’ madness was attributed to their unnatural diet which consisted of the flesh[14] of unsuspecting guests or strangers to the island.[15] Some versions of the myth say that the mares also expelled fire when they breathed.[16] The Mares, which were the terror of Thrace, were kept tethered by iron chains to a bronze manger in the now vanished city of Tirida[17] and were named Podargos (the swift), Lampon (the shining), Xanthos (the yellow) and Deinos (or Deinus, the terrible).[18] Although very similar, there are slight variances in the exact details regarding the mares’ capture.
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+ In one version, Heracles brought a number of volunteers to help him capture the giant horses.[17] After overpowering Diomedes’ men, Heracles broke the chains that tethered the horses and drove the mares down to sea. Unaware that the mares were man-eating and uncontrollable, Heracles left them in the charge of his favored companion, Abderus, while he left to fight Diomedes. Upon his return, Heracles found that the boy was eaten. As revenge, Heracles fed Diomedes to his own horses and then founded Abdera next to the boy's tomb.[15]
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+
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+ In another version, Heracles, who was visiting the island, stayed awake so that he didn't have his throat cut by Diomedes in the night, and cut the chains binding the horses once everyone was asleep. Having scared the horses onto the high ground of a knoll, Heracles quickly dug a trench through the peninsula, filling it with water and thus flooding the low-lying plain. When Diomedes and his men turned to flee, Heracles killed them with an axe (or a club[17]), and fed Diomedes’ body to the horses to calm them.
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+ In yet another version, Heracles first captured Diomedes and fed him to the mares before releasing them. Only after realizing that their King was dead did his men, the Bistonians,[15][17] attack Heracles. Upon seeing the mares charging at them, led in a chariot by Abderus, the Bistonians turned and fled.
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+ All versions have eating human flesh make the horses calmer, giving Heracles the opportunity to bind their mouths shut, and easily take them back to King Eurystheus, who dedicated the horses to Hera.[19] In some versions, they were allowed to roam freely around Argos, having become permanently calm, but in others, Eurystheus ordered the horses taken to Olympus to be sacrificed to Zeus, but Zeus refused them, and sent wolves, lions, and bears to kill them.[20] Roger Lancelyn Green states in his Tales of the Greek Heroes that the mares’ descendants were used in the Trojan War, and survived even to the time of Alexander the Great.[17][21] After the incident, Eurystheus sent Heracles to bring back Hippolyta's Girdle.
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+ Eurystheus' daughter Admete wanted the Belt of Hippolyta, queen of the Amazons, a gift from her father Ares. To please his daughter, Eurystheus ordered Heracles to retrieve the belt as his ninth labour.
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+ Taking a band of friends with him, Heracles set sail, stopping at the island of Paros, which was inhabited by some sons of Minos. The sons killed two of Heracles' companions, an act which set Heracles on a rampage. He killed two of the sons of Minos and threatened the other inhabitants until he was offered two men to replace his fallen companions. Heracles agreed and took two of Minos' grandsons, Alcaeus and Sthenelus. They continued their voyage and landed at the court of Lycus, whom Heracles defended in a battle against King Mygdon of Bebryces. After killing King Mygdon, Heracles gave much of the land to his friend Lycus. Lycus called the land Heraclea. The crew then set off for Themiscyra, where Hippolyta lived.
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+ All would have gone well for Heracles had it not been for Hera. Hippolyta, impressed with Heracles and his exploits, agreed to give him the belt and would have done so had Hera not disguised herself and walked among the Amazons sowing seeds of distrust. She claimed the strangers were plotting to carry off the queen of the Amazons. Alarmed, the women set off on horseback to confront Heracles. When Heracles saw them, he thought Hippolyta had been plotting such treachery all along and had never meant to hand over the belt, so he killed her, took the belt and returned to Eurystheus.
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+ The tenth labour was to obtain the Cattle of the three-bodied giant Geryon. In the fullest account in the Bibliotheca of Pseudo-Apollodorus,[22] Heracles had to go to the island of Erytheia in the far west (sometimes identified with the Hesperides, or with the island which forms the city of Cádiz) to get the cattle. On the way there, he crossed the Libyan desert[23] and became so frustrated at the heat that he shot an arrow at the Sun. The sun-god Helios "in admiration of his courage" gave Heracles the golden cup Helios used to sail across the sea from west to east each night. Heracles rode the cup to Erytheia; Heracles in the cup was a favorite motif on black-figure pottery.[citation needed] Such a magical conveyance undercuts any literal geography for Erytheia, the "red island" of the sunset.
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+ When Heracles landed at Erytheia, he was confronted by the two-headed dog Orthrus. With one blow from his olive-wood club, Heracles killed Orthrus. Eurytion the herdsman came to assist Orthrus, but Heracles dealt with him the same way.
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+ On hearing the commotion, Geryon sprang into action, carrying three shields and three spears, and wearing three helmets. He attacked Heracles at the River Anthemus, but was slain by one of Heracles' poisoned arrows. Heracles shot so forcefully that the arrow pierced Geryon's forehead, "and Geryon bent his neck over to one side, like a poppy that spoils its delicate shapes, shedding its petals all at once."[24]
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+ Heracles then had to herd the cattle back to Eurystheus. In Roman versions of the narrative, Heracles drove the cattle over the Aventine Hill on the future site of Rome. The giant Cacus, who lived there, stole some of the cattle as Heracles slept, making the cattle walk backwards so that they left no trail, a repetition of the trick of the young Hermes. According to some versions, Heracles drove his remaining cattle past the cave, where Cacus had hidden the stolen animals, and they began calling out to each other. In other versions, Cacus' sister Caca told Heracles where he was. Heracles then killed Cacus, and set up an altar on the spot, later the site of Rome's Forum Boarium (the cattle market).
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+ To annoy Heracles, Hera sent a gadfly to bite the cattle, irritate them, and scatter them. Within a year, Heracles retrieved them. Hera then sent a flood which raised the level of a river so much that Heracles could not cross with the cattle. He piled stones into the river to make the water shallower. When he finally reached the court of Eurystheus, the cattle were sacrificed to Hera.
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+ After Heracles completed the first ten labours, Eurystheus gave him two more, claiming that slaying the Hydra did not count (because Iolaus helped Heracles), neither did cleaning the Augean Stables (either because he was paid for the job or because the rivers did the work).
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+ The first additional labour was to steal three of the golden apples from the garden of the Hesperides. Heracles first caught the Old Man of the Sea, the shapeshifting sea god,[25] to learn where the Garden of the Hesperides was located.[26]
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+ In some variations, Heracles, either at the start or at the end of this task, meets Antaeus, who was invincible as long as he touched his mother, Gaia, the Earth. Heracles killed Antaeus by holding him aloft and crushing him in a bear hug.[27]
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+ Herodotus claims that Heracles stopped in Egypt, where King Busiris decided to make him the yearly sacrifice, but Heracles burst out of his chains.
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+ Heracles finally made his way to the garden of the Hesperides, where he encountered Atlas holding up the heavens on his shoulders. Heracles persuaded Atlas to get the three golden Apples for him by offering to hold up the heavens in his place for a little while. Atlas could get the apples because, in this version, he was the father or otherwise related to the Hesperides. This would have made the labour – like the Hydra and the Augean stables – void because Heracles had received help. When Atlas returned, he decided that he did not want to take the heavens back, and instead offered to deliver the apples himself, but Heracles tricked him by agreeing to remain in place of Atlas on the condition that Atlas relieve him temporarily while Heracles adjusted his cloak. Atlas agreed, but Heracles reneged and walked away with the apples. According to an alternative version, Heracles slew Ladon, the dragon who guarded the apples instead. Eurystheus was furious that Heracles had accomplished something that Eurystheus thought could not possibly be done.
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+ The twelfth and final labour was the capture of Cerberus, the three-headed, dragon-tailed dog that was the guardian of the gates of the Underworld. To prepare for his descent into the Underworld, Heracles went to Eleusis (or Athens) to be initiated in the Eleusinian Mysteries. He entered the Underworld, and Hermes and Athena were his guides.
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+ While in the Underworld, Heracles met Theseus and Pirithous. The two companions had been imprisoned by Hades for attempting to kidnap Persephone. One tradition tells of snakes coiling around their legs, then turning into stone; another that Hades feigned hospitality and prepared a feast inviting them to sit. They unknowingly sat in chairs of forgetfulness and were permanently ensnared. When Heracles had pulled Theseus first from his chair, some of his thigh stuck to it (this explains the supposedly lean thighs of Athenians), but the Earth shook at the attempt to liberate Pirithous, whose desire to have the goddess for himself was so insulting he was doomed to stay behind.
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+ Heracles found Hades and asked permission to bring Cerberus to the surface, which Hades agreed to if Heracles could subdue the beast without using weapons. Heracles overpowered Cerberus with his bare hands and slung the beast over his back. He carried Cerberus out of the Underworld through a cavern entrance in the Peloponnese and brought it to Eurystheus, who again fled into his pithos. Eurystheus begged Heracles to return Cerberus to the Underworld, offering in return to release him from any further labours when Cerberus disappeared back to his master.
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+ After completing the Twelve Labours, one tradition says Heracles joined Jason and the Argonauts in their quest for the Golden Fleece. However, Herodorus (c. 400 BC) disputed this and denied Heracles ever sailed with the Argonauts. A separate tradition (e.g. Argonautica) has Heracles accompany the Argonauts, but he did not travel with them as far as Colchis.
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+ Some ancient Greeks found allegorical meanings of a moral, psychological or philosophical nature in the Labours of Heracles. This trend became more prominent in the Renaissance.[28] For example, Heraclitus the Grammarian wrote in his Homeric Problems:
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+ I turn to Heracles. We must not suppose he attained such power in those days as a result of his physical strength. Rather, he was a man of intellect, an initiate in heavenly wisdom, who, as it were, shed light on philosophy, which had been hidden in deep darkness. The most authoritative of the Stoics agree with this account.... The (Erymanthian) boar which he overcame is the common incontinence of men; the (Nemean) lion is the indiscriminate rush towards improper goals; in the same way, by fettering irrational passions he gave rise to the belief that he had fettered the violent (Cretan) bull. He banished cowardice also from the world, in the shape of the hind of Ceryneia. There was another "labor" too, not properly so called, in which he cleared out the mass of dung (from the Augean stables) — in other words, the foulness that disfigures humanity. The (Stymphalian) birds he scattered are the windy hopes that feed our lives; the many-headed hydra that he burned, as it were, with the fires of exhortation, is pleasure, which begins to grow again as soon as it is cut out.
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1
+ The Twelve Labours of Heracles or Hercules (Greek: οἱ Ἡρακλέους ἆθλοι, hoi Hērakléous âthloi)[1][2] are a series of episodes concerning a penance carried out by Heracles, the greatest of the Greek heroes, whose name was later romanised as Hercules. They were accomplished at the service of King Eurystheus. The episodes were later connected by a continuous narrative. The establishment of a fixed cycle of twelve labours was attributed by the Greeks to an epic poem, now lost, written by Peisander, dated about 600 BC.[3] After Heracles killed his wife and children, he went to the oracle at Delphi. He prayed to the god Apollo for guidance. Heracles was told to serve the king of Mycenae, Eurystheus, for ten years. During this time, he is sent to perform a series of difficult feats, called labours.[4]
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+ Driven mad by Hera (queen of the gods), Heracles slew his sons by his wife Megara.[5] After recovering his sanity, Heracles deeply regretted his actions; he was purified by King Thespius, then traveled to Delphi to inquire how he could atone for his actions. Pythia, the Oracle of Delphi, advised him to go to Tiryns and serve his cousin, King Eurystheus, for ten years, performing whatever labours Eurystheus might set him; in return, he would be rewarded with immortality. Heracles despaired at this, loathing to serve a man whom he knew to be far inferior to himself, yet fearing to oppose his father Zeus. Eventually, he placed himself at Eurystheus's disposal.
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+ Eurystheus originally ordered Heracles to perform ten labours. Heracles accomplished these tasks, but Eurystheus refused to recognize two: the slaying of the Lernaean Hydra, as Heracles' nephew and charioteer Iolaus had helped him; and the cleansing of the Augeas, because Heracles accepted payment for the labour. Eurystheus set two more tasks (fetching the Golden Apples of Hesperides and capturing Cerberus), which Heracles also performed, bringing the total number of tasks to twelve.
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+ As they survive, the labours of Heracles are not recounted in any single place, but must be reassembled from many sources. Ruck and Staples[6] assert that there is no one way to interpret the labours, but that six were located in the Peloponnese, culminating with the rededication of Olympia. Six others took the hero farther afield, to places that were, per Ruck, "all previously strongholds of Hera or the 'Goddess' and were Entrances to the Netherworld".[6] In each case, the pattern was the same: Heracles was sent to kill or subdue, or to fetch back for Eurystheus (as Hera's representative) a magical animal or plant.
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+ A famous depiction of the labours in Greek sculpture is found on the metopes of the Temple of Zeus at Olympia, which date to the 450s BC.[citation needed]
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+ In his labours, Heracles was sometimes accompanied by a male companion (an eromenos), according to Licymnius[citation needed] and others, such as Iolaus, his nephew. Although he was supposed to perform only ten labours, this assistance led to two labours being disqualified: Eurystheus refused to recognize slaying the Hydra, because Iolaus helped him, and the cleansing of the Augean stables, because Heracles was paid for his services and/or because the rivers did the work. Several of the labours involved the offspring (by various accounts) of Typhon and his mate Echidna, all overcome by Heracles.
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+ A traditional order of the labours found in the Bibliotheca[7] by Pseudo-Apollodorus is:
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+ Heracles wandered the area until he came to the town of Cleonae. There he met a boy who said that if Heracles slew the Nemean lion and returned alive within 30 days, the town would sacrifice a lion to Zeus, but if he did not return within 30 days or if he died, the boy would sacrifice himself to Zeus. Another version claims that he met Molorchos, a shepherd who had lost his son to the lion, saying that if he came back within 30 days, a ram would be sacrificed to Zeus. If he did not return within 30 days, it would be sacrificed to the dead Heracles as a mourning offering.
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+ While searching for the lion, Heracles fletched some arrows to use against it, not knowing that its golden fur was impenetrable. When he found and shot the lion, firing at it with his bow, he discovered the fur's protective property as the arrow bounced harmlessly off the creature's thigh. After some time, Heracles made the lion return to his cave. The cave had two entrances, one of which Heracles blocked; he then entered the other. In those dark and close quarters, Heracles stunned the beast with his club and, using his immense strength, strangled it to death. During the fight the lion bit off one of his fingers.[8] Others say that he shot arrows at it, eventually shooting it in the unarmored mouth. After slaying the lion, he tried to skin it with a knife from his belt, but failed. He then tried sharpening the knife with a stone and even tried with the stone itself. Finally, Athena, noticing the hero's plight, told Heracles to use one of the lion's own claws to skin the pelt. Others say that Heracles' armor was, in fact, the hide of the Lion of Cithaeron.
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+ When he returned on the 30th day carrying the carcass of the lion on his shoulders, King Eurystheus was amazed and terrified. Eurystheus forbade him ever again to enter the city; from then on he was to display the fruits of his labours outside the city gates. Eurystheus would then tell Heracles his tasks through a herald, not personally. Eurystheus even had a large bronze jar made for him in which to hide from Heracles if need be. Eurystheus then warned him that the tasks would become increasingly difficult.
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+ Heracles' second labour was to slay the Lernaean Hydra, which Hera had raised just to slay Heracles. Upon reaching the swamp near Lake Lerna, where the Hydra dwelt, Heracles used a cloth to cover his mouth and nose to protect himself from the poisonous fumes. He fired flaming arrows into the Hydra's lair, the spring of Amymone, a deep cave that it only came out of to terrorize neighboring villages.[9] He then confronted the Hydra, wielding a harvesting sickle (according to some early vase-paintings), a sword or his famed club. Ruck and Staples (1994: 170) have pointed out that the chthonic creature's reaction was botanical: upon cutting off each of its heads he found that two grew back, an expression of the hopelessness of such a struggle for any but the hero. Additionally, one of the Hydra's heads - the middle one - was immortal.
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+ The details of the struggle are explicit in the Bibliotheca (2.5.2): realizing that he could not defeat the Hydra in this way, Heracles called on his nephew Iolaus for help. His nephew then came upon the idea (possibly inspired by Athena) of using a firebrand to scorch the neck stumps after each decapitation. Heracles cut off each head and Iolaus cauterized the open stumps. Seeing that Heracles was winning the struggle, Hera sent a giant crab to distract him. He crushed it under his mighty foot. He cut off the Hydra's one immortal head with a golden sword given to him by Athena. Heracles placed it under a great rock on the sacred way between Lerna and Elaius (Kerenyi 1959:144), and dipped his arrows in the Hydra's poisonous blood, and so his second task was complete. The alternative version of this myth is that after cutting off one head, he then dipped his sword in it and used its venom to burn each head so it could not grow back. Hera, upset that Heracles had slain the beast she raised to kill him, placed it in the dark blue vault of the sky as the constellation Hydra. She then turned the crab into the constellation Cancer.
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+ Later, Heracles used an arrow dipped in the Hydra's poisonous blood to kill the centaur Nessus; and Nessus's tainted blood was applied to the Tunic of Nessus, by which the centaur had his posthumous revenge. Both Strabo and Pausanias report that the stench of the river Anigrus in Elis, making all the fish of the river inedible, was reputed to be due to the Hydra's venom, washed from the arrows Heracles used on the centaur.[10]
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+ Eurystheus and Hera were greatly angered that Heracles had survived the Nemean Lion and the Lernaean Hydra. For the third labour, they found a task which they thought would spell doom for the hero. It was not slaying a beast or monster, as it had already been established that Heracles could overcome even the most fearsome opponents. Instead, Eurystheus ordered him to capture the Ceryneian Hind, which was so fast that it could outrun an arrow.
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+ After beginning the search, Heracles awoke from sleeping and saw the hind by the glint on its antlers. Heracles then chased the hind on foot for a full year through Greece, Thrace, Istria, and the land of the Hyperboreans. In some versions, he captured the hind while it slept, rendering it lame with a trap net. In other versions, he encountered Artemis in her temple; she told him to leave the hind and tell Eurystheus all that had happened, and his third labour would be considered to be completed. Yet another version claims that Heracles trapped the Hind with an arrow between its forelegs.
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+ Eurystheus had given Heracles this task hoping to incite Artemis' anger at Heracles for his desecration of her sacred animal. As he was returning with the hind, Heracles encountered Artemis and her brother Apollo. He begged the goddess for forgiveness, explaining that he had to catch it as part of his penance, but he promised to return it. Artemis forgave him, foiling Eurystheus' plan to have her punish him.
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+ Upon bringing the hind to Eurystheus, he was told that it was to become part of the King's menagerie. Heracles knew that he had to return the hind as he had promised, so he agreed to hand it over on the condition that Eurystheus himself come out and take it from him. The King came out, but the moment that Heracles let the hind go, it sprinted back to its mistress and Heracles left, saying that Eurystheus had not been quick enough.
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+ Eurystheus was disappointed that Heracles had overcome yet another creature and was humiliated by the hind's escape, so he assigned Heracles another dangerous task. By some accounts, the fourth labour was to bring the fearsome Erymanthian Boar back to Eurystheus alive (there is no single definitive telling of the labours). On the way to Mount Erymanthos where the boar lived, Heracles visited Pholus ("caveman"), a kind and hospitable centaur and old friend. Heracles ate with Pholus in his cavern (though the centaur devoured his meat raw) and asked for wine. Pholus had only one jar of wine, a gift from Dionysus to all the centaurs on Mount Erymanthos. Heracles convinced him to open it, and the smell attracted the other centaurs. They did not understand that wine needs to be tempered with water, became drunk, and attacked Heracles. Heracles shot at them with his poisonous arrows, killing many, and the centaurs retreated all the way to Chiron's cave.
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+ Pholus was curious why the arrows caused so much death. He picked one up but dropped it, and the arrow stabbed his hoof, poisoning him. One version states that a stray arrow hit Chiron as well. He was immortal, but he still felt the pain. Chiron's pain was so great that he volunteered to give up his immortality and take the place of Prometheus, who had been chained to the top of a mountain to have his liver eaten daily by an eagle. Prometheus' torturer, the eagle, continued its torture on Chiron, so Heracles shot it dead with an arrow. It is generally accepted that the tale was meant to show Heracles as being the recipient of Chiron's surrendered immortality. However, this tale contradicts the fact that Chiron later taught Achilles. The tale of the centaurs sometimes appears in other parts of the twelve labours, as does the freeing of Prometheus.
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+ Heracles had visited Chiron to gain advice on how to catch the boar, and Chiron had told him to drive it into thick snow, which sets this labour in mid-winter. Heracles caught the boar, bound it, and carried it back to Eurystheus, who was frightened of it and ducked down in his half-buried storage pithos, begging Heracles to get rid of the beast.
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+ The fifth labour was to clean the stables of King Augeas. This assignment was intended to be both humiliating (rather than impressive, as the previous labours had been) and impossible, since the livestock were divinely healthy (and immortal) and therefore produced an enormous quantity of dung. The Augean Stables (/ɔːˈdʒiːən/) had not been cleaned in over 30 years, and over 1,000 cattle lived there. However, Heracles succeeded by re-routing the rivers Alpheus and Peneus to wash out the filth.
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+ Before starting on the task, Heracles had asked Augeas for one-tenth of the cattle if he finished the task in one day, and Augeas agreed. But afterwards Augeas refused to honour the agreement on the grounds that Heracles had been ordered to carry out the task by Eurystheus anyway. Heracles claimed his reward in court, and was supported by Augeas' son Phyleus. Augeas banished them both before the court had ruled. Heracles returned, slew Augeas, and gave his kingdom to Phyleus. Heracles then founded the Olympic Games.
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+ The success of this labour was ultimately discounted as the rushing waters had done the work of cleaning the stables and because Heracles was paid for doing the labour.
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+ Eurystheus said that Heracles still had seven labours to perform.[11]
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+
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+ The sixth labour was to defeat the Stymphalian birds, man-eating birds with beaks made of bronze and sharp metallic feathers they could launch at their victims. They were sacred to Ares, the god of war. Furthermore, their dung was highly toxic. They had migrated to Lake Stymphalia in Arcadia, where they bred quickly and took over the countryside, destroying local crops, fruit trees, and townspeople. Heracles could not go too far into the swamp, for it would not support his weight. Athena, noticing the hero's plight, gave Heracles a rattle which Hephaestus had made especially for the occasion. Heracles shook the rattle and frightened the birds into the air. Heracles then shot many of them with his arrows. The rest flew far away, never to return. The Argonauts would later encounter them.
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+
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+ The seventh labour was to capture the Cretan Bull, father of the Minotaur. Heracles sailed to Crete, where King Minos gave Heracles permission to take the bull away and even offered him assistance (which Heracles declined, plausibly because he did not want the labour to be discounted as before).[12] The bull had been wreaking havoc on Crete by uprooting crops and leveling orchard walls. Heracles sneaked up behind the bull and then used his hands to throttle it (stopping before it was killed), and then shipped it back to Tiryns. Eurystheus, who hid in his pithos at first sight of the creature, wanted to sacrifice the bull to Hera, who hated Heracles. She refused the sacrifice because it reflected glory on Heracles. The bull was released and wandered into Marathon, becoming known as the Marathonian Bull.[12] Theseus would later sacrifice the bull to Athena and/or Apollo.
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+
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+ As the eighth of his Twelve Labours, also categorised as the second of the Non-Peloponneisan labours,[13] Heracles was sent by King Eurystheus to steal the Mares from Diomedes. The mares’ madness was attributed to their unnatural diet which consisted of the flesh[14] of unsuspecting guests or strangers to the island.[15] Some versions of the myth say that the mares also expelled fire when they breathed.[16] The Mares, which were the terror of Thrace, were kept tethered by iron chains to a bronze manger in the now vanished city of Tirida[17] and were named Podargos (the swift), Lampon (the shining), Xanthos (the yellow) and Deinos (or Deinus, the terrible).[18] Although very similar, there are slight variances in the exact details regarding the mares’ capture.
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+ In one version, Heracles brought a number of volunteers to help him capture the giant horses.[17] After overpowering Diomedes’ men, Heracles broke the chains that tethered the horses and drove the mares down to sea. Unaware that the mares were man-eating and uncontrollable, Heracles left them in the charge of his favored companion, Abderus, while he left to fight Diomedes. Upon his return, Heracles found that the boy was eaten. As revenge, Heracles fed Diomedes to his own horses and then founded Abdera next to the boy's tomb.[15]
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+
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+ In another version, Heracles, who was visiting the island, stayed awake so that he didn't have his throat cut by Diomedes in the night, and cut the chains binding the horses once everyone was asleep. Having scared the horses onto the high ground of a knoll, Heracles quickly dug a trench through the peninsula, filling it with water and thus flooding the low-lying plain. When Diomedes and his men turned to flee, Heracles killed them with an axe (or a club[17]), and fed Diomedes’ body to the horses to calm them.
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+ In yet another version, Heracles first captured Diomedes and fed him to the mares before releasing them. Only after realizing that their King was dead did his men, the Bistonians,[15][17] attack Heracles. Upon seeing the mares charging at them, led in a chariot by Abderus, the Bistonians turned and fled.
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+ All versions have eating human flesh make the horses calmer, giving Heracles the opportunity to bind their mouths shut, and easily take them back to King Eurystheus, who dedicated the horses to Hera.[19] In some versions, they were allowed to roam freely around Argos, having become permanently calm, but in others, Eurystheus ordered the horses taken to Olympus to be sacrificed to Zeus, but Zeus refused them, and sent wolves, lions, and bears to kill them.[20] Roger Lancelyn Green states in his Tales of the Greek Heroes that the mares’ descendants were used in the Trojan War, and survived even to the time of Alexander the Great.[17][21] After the incident, Eurystheus sent Heracles to bring back Hippolyta's Girdle.
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+
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+ Eurystheus' daughter Admete wanted the Belt of Hippolyta, queen of the Amazons, a gift from her father Ares. To please his daughter, Eurystheus ordered Heracles to retrieve the belt as his ninth labour.
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+ Taking a band of friends with him, Heracles set sail, stopping at the island of Paros, which was inhabited by some sons of Minos. The sons killed two of Heracles' companions, an act which set Heracles on a rampage. He killed two of the sons of Minos and threatened the other inhabitants until he was offered two men to replace his fallen companions. Heracles agreed and took two of Minos' grandsons, Alcaeus and Sthenelus. They continued their voyage and landed at the court of Lycus, whom Heracles defended in a battle against King Mygdon of Bebryces. After killing King Mygdon, Heracles gave much of the land to his friend Lycus. Lycus called the land Heraclea. The crew then set off for Themiscyra, where Hippolyta lived.
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+ All would have gone well for Heracles had it not been for Hera. Hippolyta, impressed with Heracles and his exploits, agreed to give him the belt and would have done so had Hera not disguised herself and walked among the Amazons sowing seeds of distrust. She claimed the strangers were plotting to carry off the queen of the Amazons. Alarmed, the women set off on horseback to confront Heracles. When Heracles saw them, he thought Hippolyta had been plotting such treachery all along and had never meant to hand over the belt, so he killed her, took the belt and returned to Eurystheus.
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+ The tenth labour was to obtain the Cattle of the three-bodied giant Geryon. In the fullest account in the Bibliotheca of Pseudo-Apollodorus,[22] Heracles had to go to the island of Erytheia in the far west (sometimes identified with the Hesperides, or with the island which forms the city of Cádiz) to get the cattle. On the way there, he crossed the Libyan desert[23] and became so frustrated at the heat that he shot an arrow at the Sun. The sun-god Helios "in admiration of his courage" gave Heracles the golden cup Helios used to sail across the sea from west to east each night. Heracles rode the cup to Erytheia; Heracles in the cup was a favorite motif on black-figure pottery.[citation needed] Such a magical conveyance undercuts any literal geography for Erytheia, the "red island" of the sunset.
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+ When Heracles landed at Erytheia, he was confronted by the two-headed dog Orthrus. With one blow from his olive-wood club, Heracles killed Orthrus. Eurytion the herdsman came to assist Orthrus, but Heracles dealt with him the same way.
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+ On hearing the commotion, Geryon sprang into action, carrying three shields and three spears, and wearing three helmets. He attacked Heracles at the River Anthemus, but was slain by one of Heracles' poisoned arrows. Heracles shot so forcefully that the arrow pierced Geryon's forehead, "and Geryon bent his neck over to one side, like a poppy that spoils its delicate shapes, shedding its petals all at once."[24]
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+ Heracles then had to herd the cattle back to Eurystheus. In Roman versions of the narrative, Heracles drove the cattle over the Aventine Hill on the future site of Rome. The giant Cacus, who lived there, stole some of the cattle as Heracles slept, making the cattle walk backwards so that they left no trail, a repetition of the trick of the young Hermes. According to some versions, Heracles drove his remaining cattle past the cave, where Cacus had hidden the stolen animals, and they began calling out to each other. In other versions, Cacus' sister Caca told Heracles where he was. Heracles then killed Cacus, and set up an altar on the spot, later the site of Rome's Forum Boarium (the cattle market).
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+ To annoy Heracles, Hera sent a gadfly to bite the cattle, irritate them, and scatter them. Within a year, Heracles retrieved them. Hera then sent a flood which raised the level of a river so much that Heracles could not cross with the cattle. He piled stones into the river to make the water shallower. When he finally reached the court of Eurystheus, the cattle were sacrificed to Hera.
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+ After Heracles completed the first ten labours, Eurystheus gave him two more, claiming that slaying the Hydra did not count (because Iolaus helped Heracles), neither did cleaning the Augean Stables (either because he was paid for the job or because the rivers did the work).
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+ The first additional labour was to steal three of the golden apples from the garden of the Hesperides. Heracles first caught the Old Man of the Sea, the shapeshifting sea god,[25] to learn where the Garden of the Hesperides was located.[26]
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+ In some variations, Heracles, either at the start or at the end of this task, meets Antaeus, who was invincible as long as he touched his mother, Gaia, the Earth. Heracles killed Antaeus by holding him aloft and crushing him in a bear hug.[27]
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+ Herodotus claims that Heracles stopped in Egypt, where King Busiris decided to make him the yearly sacrifice, but Heracles burst out of his chains.
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+ Heracles finally made his way to the garden of the Hesperides, where he encountered Atlas holding up the heavens on his shoulders. Heracles persuaded Atlas to get the three golden Apples for him by offering to hold up the heavens in his place for a little while. Atlas could get the apples because, in this version, he was the father or otherwise related to the Hesperides. This would have made the labour – like the Hydra and the Augean stables – void because Heracles had received help. When Atlas returned, he decided that he did not want to take the heavens back, and instead offered to deliver the apples himself, but Heracles tricked him by agreeing to remain in place of Atlas on the condition that Atlas relieve him temporarily while Heracles adjusted his cloak. Atlas agreed, but Heracles reneged and walked away with the apples. According to an alternative version, Heracles slew Ladon, the dragon who guarded the apples instead. Eurystheus was furious that Heracles had accomplished something that Eurystheus thought could not possibly be done.
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+ The twelfth and final labour was the capture of Cerberus, the three-headed, dragon-tailed dog that was the guardian of the gates of the Underworld. To prepare for his descent into the Underworld, Heracles went to Eleusis (or Athens) to be initiated in the Eleusinian Mysteries. He entered the Underworld, and Hermes and Athena were his guides.
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+ While in the Underworld, Heracles met Theseus and Pirithous. The two companions had been imprisoned by Hades for attempting to kidnap Persephone. One tradition tells of snakes coiling around their legs, then turning into stone; another that Hades feigned hospitality and prepared a feast inviting them to sit. They unknowingly sat in chairs of forgetfulness and were permanently ensnared. When Heracles had pulled Theseus first from his chair, some of his thigh stuck to it (this explains the supposedly lean thighs of Athenians), but the Earth shook at the attempt to liberate Pirithous, whose desire to have the goddess for himself was so insulting he was doomed to stay behind.
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+ Heracles found Hades and asked permission to bring Cerberus to the surface, which Hades agreed to if Heracles could subdue the beast without using weapons. Heracles overpowered Cerberus with his bare hands and slung the beast over his back. He carried Cerberus out of the Underworld through a cavern entrance in the Peloponnese and brought it to Eurystheus, who again fled into his pithos. Eurystheus begged Heracles to return Cerberus to the Underworld, offering in return to release him from any further labours when Cerberus disappeared back to his master.
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+ After completing the Twelve Labours, one tradition says Heracles joined Jason and the Argonauts in their quest for the Golden Fleece. However, Herodorus (c. 400 BC) disputed this and denied Heracles ever sailed with the Argonauts. A separate tradition (e.g. Argonautica) has Heracles accompany the Argonauts, but he did not travel with them as far as Colchis.
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+ Some ancient Greeks found allegorical meanings of a moral, psychological or philosophical nature in the Labours of Heracles. This trend became more prominent in the Renaissance.[28] For example, Heraclitus the Grammarian wrote in his Homeric Problems:
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+ I turn to Heracles. We must not suppose he attained such power in those days as a result of his physical strength. Rather, he was a man of intellect, an initiate in heavenly wisdom, who, as it were, shed light on philosophy, which had been hidden in deep darkness. The most authoritative of the Stoics agree with this account.... The (Erymanthian) boar which he overcame is the common incontinence of men; the (Nemean) lion is the indiscriminate rush towards improper goals; in the same way, by fettering irrational passions he gave rise to the belief that he had fettered the violent (Cretan) bull. He banished cowardice also from the world, in the shape of the hind of Ceryneia. There was another "labor" too, not properly so called, in which he cleared out the mass of dung (from the Augean stables) — in other words, the foulness that disfigures humanity. The (Stymphalian) birds he scattered are the windy hopes that feed our lives; the many-headed hydra that he burned, as it were, with the fires of exhortation, is pleasure, which begins to grow again as soon as it is cut out.
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1
+ The Twelve Labours of Heracles or Hercules (Greek: οἱ Ἡρακλέους ἆθλοι, hoi Hērakléous âthloi)[1][2] are a series of episodes concerning a penance carried out by Heracles, the greatest of the Greek heroes, whose name was later romanised as Hercules. They were accomplished at the service of King Eurystheus. The episodes were later connected by a continuous narrative. The establishment of a fixed cycle of twelve labours was attributed by the Greeks to an epic poem, now lost, written by Peisander, dated about 600 BC.[3] After Heracles killed his wife and children, he went to the oracle at Delphi. He prayed to the god Apollo for guidance. Heracles was told to serve the king of Mycenae, Eurystheus, for ten years. During this time, he is sent to perform a series of difficult feats, called labours.[4]
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+ Driven mad by Hera (queen of the gods), Heracles slew his sons by his wife Megara.[5] After recovering his sanity, Heracles deeply regretted his actions; he was purified by King Thespius, then traveled to Delphi to inquire how he could atone for his actions. Pythia, the Oracle of Delphi, advised him to go to Tiryns and serve his cousin, King Eurystheus, for ten years, performing whatever labours Eurystheus might set him; in return, he would be rewarded with immortality. Heracles despaired at this, loathing to serve a man whom he knew to be far inferior to himself, yet fearing to oppose his father Zeus. Eventually, he placed himself at Eurystheus's disposal.
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+ Eurystheus originally ordered Heracles to perform ten labours. Heracles accomplished these tasks, but Eurystheus refused to recognize two: the slaying of the Lernaean Hydra, as Heracles' nephew and charioteer Iolaus had helped him; and the cleansing of the Augeas, because Heracles accepted payment for the labour. Eurystheus set two more tasks (fetching the Golden Apples of Hesperides and capturing Cerberus), which Heracles also performed, bringing the total number of tasks to twelve.
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+ As they survive, the labours of Heracles are not recounted in any single place, but must be reassembled from many sources. Ruck and Staples[6] assert that there is no one way to interpret the labours, but that six were located in the Peloponnese, culminating with the rededication of Olympia. Six others took the hero farther afield, to places that were, per Ruck, "all previously strongholds of Hera or the 'Goddess' and were Entrances to the Netherworld".[6] In each case, the pattern was the same: Heracles was sent to kill or subdue, or to fetch back for Eurystheus (as Hera's representative) a magical animal or plant.
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+ A famous depiction of the labours in Greek sculpture is found on the metopes of the Temple of Zeus at Olympia, which date to the 450s BC.[citation needed]
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+ In his labours, Heracles was sometimes accompanied by a male companion (an eromenos), according to Licymnius[citation needed] and others, such as Iolaus, his nephew. Although he was supposed to perform only ten labours, this assistance led to two labours being disqualified: Eurystheus refused to recognize slaying the Hydra, because Iolaus helped him, and the cleansing of the Augean stables, because Heracles was paid for his services and/or because the rivers did the work. Several of the labours involved the offspring (by various accounts) of Typhon and his mate Echidna, all overcome by Heracles.
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+ A traditional order of the labours found in the Bibliotheca[7] by Pseudo-Apollodorus is:
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+ Heracles wandered the area until he came to the town of Cleonae. There he met a boy who said that if Heracles slew the Nemean lion and returned alive within 30 days, the town would sacrifice a lion to Zeus, but if he did not return within 30 days or if he died, the boy would sacrifice himself to Zeus. Another version claims that he met Molorchos, a shepherd who had lost his son to the lion, saying that if he came back within 30 days, a ram would be sacrificed to Zeus. If he did not return within 30 days, it would be sacrificed to the dead Heracles as a mourning offering.
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+ While searching for the lion, Heracles fletched some arrows to use against it, not knowing that its golden fur was impenetrable. When he found and shot the lion, firing at it with his bow, he discovered the fur's protective property as the arrow bounced harmlessly off the creature's thigh. After some time, Heracles made the lion return to his cave. The cave had two entrances, one of which Heracles blocked; he then entered the other. In those dark and close quarters, Heracles stunned the beast with his club and, using his immense strength, strangled it to death. During the fight the lion bit off one of his fingers.[8] Others say that he shot arrows at it, eventually shooting it in the unarmored mouth. After slaying the lion, he tried to skin it with a knife from his belt, but failed. He then tried sharpening the knife with a stone and even tried with the stone itself. Finally, Athena, noticing the hero's plight, told Heracles to use one of the lion's own claws to skin the pelt. Others say that Heracles' armor was, in fact, the hide of the Lion of Cithaeron.
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+ When he returned on the 30th day carrying the carcass of the lion on his shoulders, King Eurystheus was amazed and terrified. Eurystheus forbade him ever again to enter the city; from then on he was to display the fruits of his labours outside the city gates. Eurystheus would then tell Heracles his tasks through a herald, not personally. Eurystheus even had a large bronze jar made for him in which to hide from Heracles if need be. Eurystheus then warned him that the tasks would become increasingly difficult.
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+ Heracles' second labour was to slay the Lernaean Hydra, which Hera had raised just to slay Heracles. Upon reaching the swamp near Lake Lerna, where the Hydra dwelt, Heracles used a cloth to cover his mouth and nose to protect himself from the poisonous fumes. He fired flaming arrows into the Hydra's lair, the spring of Amymone, a deep cave that it only came out of to terrorize neighboring villages.[9] He then confronted the Hydra, wielding a harvesting sickle (according to some early vase-paintings), a sword or his famed club. Ruck and Staples (1994: 170) have pointed out that the chthonic creature's reaction was botanical: upon cutting off each of its heads he found that two grew back, an expression of the hopelessness of such a struggle for any but the hero. Additionally, one of the Hydra's heads - the middle one - was immortal.
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+ The details of the struggle are explicit in the Bibliotheca (2.5.2): realizing that he could not defeat the Hydra in this way, Heracles called on his nephew Iolaus for help. His nephew then came upon the idea (possibly inspired by Athena) of using a firebrand to scorch the neck stumps after each decapitation. Heracles cut off each head and Iolaus cauterized the open stumps. Seeing that Heracles was winning the struggle, Hera sent a giant crab to distract him. He crushed it under his mighty foot. He cut off the Hydra's one immortal head with a golden sword given to him by Athena. Heracles placed it under a great rock on the sacred way between Lerna and Elaius (Kerenyi 1959:144), and dipped his arrows in the Hydra's poisonous blood, and so his second task was complete. The alternative version of this myth is that after cutting off one head, he then dipped his sword in it and used its venom to burn each head so it could not grow back. Hera, upset that Heracles had slain the beast she raised to kill him, placed it in the dark blue vault of the sky as the constellation Hydra. She then turned the crab into the constellation Cancer.
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+ Later, Heracles used an arrow dipped in the Hydra's poisonous blood to kill the centaur Nessus; and Nessus's tainted blood was applied to the Tunic of Nessus, by which the centaur had his posthumous revenge. Both Strabo and Pausanias report that the stench of the river Anigrus in Elis, making all the fish of the river inedible, was reputed to be due to the Hydra's venom, washed from the arrows Heracles used on the centaur.[10]
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+ Eurystheus and Hera were greatly angered that Heracles had survived the Nemean Lion and the Lernaean Hydra. For the third labour, they found a task which they thought would spell doom for the hero. It was not slaying a beast or monster, as it had already been established that Heracles could overcome even the most fearsome opponents. Instead, Eurystheus ordered him to capture the Ceryneian Hind, which was so fast that it could outrun an arrow.
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+ After beginning the search, Heracles awoke from sleeping and saw the hind by the glint on its antlers. Heracles then chased the hind on foot for a full year through Greece, Thrace, Istria, and the land of the Hyperboreans. In some versions, he captured the hind while it slept, rendering it lame with a trap net. In other versions, he encountered Artemis in her temple; she told him to leave the hind and tell Eurystheus all that had happened, and his third labour would be considered to be completed. Yet another version claims that Heracles trapped the Hind with an arrow between its forelegs.
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+ Eurystheus had given Heracles this task hoping to incite Artemis' anger at Heracles for his desecration of her sacred animal. As he was returning with the hind, Heracles encountered Artemis and her brother Apollo. He begged the goddess for forgiveness, explaining that he had to catch it as part of his penance, but he promised to return it. Artemis forgave him, foiling Eurystheus' plan to have her punish him.
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+ Upon bringing the hind to Eurystheus, he was told that it was to become part of the King's menagerie. Heracles knew that he had to return the hind as he had promised, so he agreed to hand it over on the condition that Eurystheus himself come out and take it from him. The King came out, but the moment that Heracles let the hind go, it sprinted back to its mistress and Heracles left, saying that Eurystheus had not been quick enough.
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+ Eurystheus was disappointed that Heracles had overcome yet another creature and was humiliated by the hind's escape, so he assigned Heracles another dangerous task. By some accounts, the fourth labour was to bring the fearsome Erymanthian Boar back to Eurystheus alive (there is no single definitive telling of the labours). On the way to Mount Erymanthos where the boar lived, Heracles visited Pholus ("caveman"), a kind and hospitable centaur and old friend. Heracles ate with Pholus in his cavern (though the centaur devoured his meat raw) and asked for wine. Pholus had only one jar of wine, a gift from Dionysus to all the centaurs on Mount Erymanthos. Heracles convinced him to open it, and the smell attracted the other centaurs. They did not understand that wine needs to be tempered with water, became drunk, and attacked Heracles. Heracles shot at them with his poisonous arrows, killing many, and the centaurs retreated all the way to Chiron's cave.
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+ Pholus was curious why the arrows caused so much death. He picked one up but dropped it, and the arrow stabbed his hoof, poisoning him. One version states that a stray arrow hit Chiron as well. He was immortal, but he still felt the pain. Chiron's pain was so great that he volunteered to give up his immortality and take the place of Prometheus, who had been chained to the top of a mountain to have his liver eaten daily by an eagle. Prometheus' torturer, the eagle, continued its torture on Chiron, so Heracles shot it dead with an arrow. It is generally accepted that the tale was meant to show Heracles as being the recipient of Chiron's surrendered immortality. However, this tale contradicts the fact that Chiron later taught Achilles. The tale of the centaurs sometimes appears in other parts of the twelve labours, as does the freeing of Prometheus.
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+ Heracles had visited Chiron to gain advice on how to catch the boar, and Chiron had told him to drive it into thick snow, which sets this labour in mid-winter. Heracles caught the boar, bound it, and carried it back to Eurystheus, who was frightened of it and ducked down in his half-buried storage pithos, begging Heracles to get rid of the beast.
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+ The fifth labour was to clean the stables of King Augeas. This assignment was intended to be both humiliating (rather than impressive, as the previous labours had been) and impossible, since the livestock were divinely healthy (and immortal) and therefore produced an enormous quantity of dung. The Augean Stables (/ɔːˈdʒiːən/) had not been cleaned in over 30 years, and over 1,000 cattle lived there. However, Heracles succeeded by re-routing the rivers Alpheus and Peneus to wash out the filth.
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+ Before starting on the task, Heracles had asked Augeas for one-tenth of the cattle if he finished the task in one day, and Augeas agreed. But afterwards Augeas refused to honour the agreement on the grounds that Heracles had been ordered to carry out the task by Eurystheus anyway. Heracles claimed his reward in court, and was supported by Augeas' son Phyleus. Augeas banished them both before the court had ruled. Heracles returned, slew Augeas, and gave his kingdom to Phyleus. Heracles then founded the Olympic Games.
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+ The success of this labour was ultimately discounted as the rushing waters had done the work of cleaning the stables and because Heracles was paid for doing the labour.
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+ Eurystheus said that Heracles still had seven labours to perform.[11]
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+ The sixth labour was to defeat the Stymphalian birds, man-eating birds with beaks made of bronze and sharp metallic feathers they could launch at their victims. They were sacred to Ares, the god of war. Furthermore, their dung was highly toxic. They had migrated to Lake Stymphalia in Arcadia, where they bred quickly and took over the countryside, destroying local crops, fruit trees, and townspeople. Heracles could not go too far into the swamp, for it would not support his weight. Athena, noticing the hero's plight, gave Heracles a rattle which Hephaestus had made especially for the occasion. Heracles shook the rattle and frightened the birds into the air. Heracles then shot many of them with his arrows. The rest flew far away, never to return. The Argonauts would later encounter them.
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+ The seventh labour was to capture the Cretan Bull, father of the Minotaur. Heracles sailed to Crete, where King Minos gave Heracles permission to take the bull away and even offered him assistance (which Heracles declined, plausibly because he did not want the labour to be discounted as before).[12] The bull had been wreaking havoc on Crete by uprooting crops and leveling orchard walls. Heracles sneaked up behind the bull and then used his hands to throttle it (stopping before it was killed), and then shipped it back to Tiryns. Eurystheus, who hid in his pithos at first sight of the creature, wanted to sacrifice the bull to Hera, who hated Heracles. She refused the sacrifice because it reflected glory on Heracles. The bull was released and wandered into Marathon, becoming known as the Marathonian Bull.[12] Theseus would later sacrifice the bull to Athena and/or Apollo.
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+
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+ As the eighth of his Twelve Labours, also categorised as the second of the Non-Peloponneisan labours,[13] Heracles was sent by King Eurystheus to steal the Mares from Diomedes. The mares’ madness was attributed to their unnatural diet which consisted of the flesh[14] of unsuspecting guests or strangers to the island.[15] Some versions of the myth say that the mares also expelled fire when they breathed.[16] The Mares, which were the terror of Thrace, were kept tethered by iron chains to a bronze manger in the now vanished city of Tirida[17] and were named Podargos (the swift), Lampon (the shining), Xanthos (the yellow) and Deinos (or Deinus, the terrible).[18] Although very similar, there are slight variances in the exact details regarding the mares’ capture.
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+ In one version, Heracles brought a number of volunteers to help him capture the giant horses.[17] After overpowering Diomedes’ men, Heracles broke the chains that tethered the horses and drove the mares down to sea. Unaware that the mares were man-eating and uncontrollable, Heracles left them in the charge of his favored companion, Abderus, while he left to fight Diomedes. Upon his return, Heracles found that the boy was eaten. As revenge, Heracles fed Diomedes to his own horses and then founded Abdera next to the boy's tomb.[15]
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+
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+ In another version, Heracles, who was visiting the island, stayed awake so that he didn't have his throat cut by Diomedes in the night, and cut the chains binding the horses once everyone was asleep. Having scared the horses onto the high ground of a knoll, Heracles quickly dug a trench through the peninsula, filling it with water and thus flooding the low-lying plain. When Diomedes and his men turned to flee, Heracles killed them with an axe (or a club[17]), and fed Diomedes’ body to the horses to calm them.
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+ In yet another version, Heracles first captured Diomedes and fed him to the mares before releasing them. Only after realizing that their King was dead did his men, the Bistonians,[15][17] attack Heracles. Upon seeing the mares charging at them, led in a chariot by Abderus, the Bistonians turned and fled.
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+ All versions have eating human flesh make the horses calmer, giving Heracles the opportunity to bind their mouths shut, and easily take them back to King Eurystheus, who dedicated the horses to Hera.[19] In some versions, they were allowed to roam freely around Argos, having become permanently calm, but in others, Eurystheus ordered the horses taken to Olympus to be sacrificed to Zeus, but Zeus refused them, and sent wolves, lions, and bears to kill them.[20] Roger Lancelyn Green states in his Tales of the Greek Heroes that the mares’ descendants were used in the Trojan War, and survived even to the time of Alexander the Great.[17][21] After the incident, Eurystheus sent Heracles to bring back Hippolyta's Girdle.
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+
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+ Eurystheus' daughter Admete wanted the Belt of Hippolyta, queen of the Amazons, a gift from her father Ares. To please his daughter, Eurystheus ordered Heracles to retrieve the belt as his ninth labour.
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+ Taking a band of friends with him, Heracles set sail, stopping at the island of Paros, which was inhabited by some sons of Minos. The sons killed two of Heracles' companions, an act which set Heracles on a rampage. He killed two of the sons of Minos and threatened the other inhabitants until he was offered two men to replace his fallen companions. Heracles agreed and took two of Minos' grandsons, Alcaeus and Sthenelus. They continued their voyage and landed at the court of Lycus, whom Heracles defended in a battle against King Mygdon of Bebryces. After killing King Mygdon, Heracles gave much of the land to his friend Lycus. Lycus called the land Heraclea. The crew then set off for Themiscyra, where Hippolyta lived.
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+ All would have gone well for Heracles had it not been for Hera. Hippolyta, impressed with Heracles and his exploits, agreed to give him the belt and would have done so had Hera not disguised herself and walked among the Amazons sowing seeds of distrust. She claimed the strangers were plotting to carry off the queen of the Amazons. Alarmed, the women set off on horseback to confront Heracles. When Heracles saw them, he thought Hippolyta had been plotting such treachery all along and had never meant to hand over the belt, so he killed her, took the belt and returned to Eurystheus.
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+ The tenth labour was to obtain the Cattle of the three-bodied giant Geryon. In the fullest account in the Bibliotheca of Pseudo-Apollodorus,[22] Heracles had to go to the island of Erytheia in the far west (sometimes identified with the Hesperides, or with the island which forms the city of Cádiz) to get the cattle. On the way there, he crossed the Libyan desert[23] and became so frustrated at the heat that he shot an arrow at the Sun. The sun-god Helios "in admiration of his courage" gave Heracles the golden cup Helios used to sail across the sea from west to east each night. Heracles rode the cup to Erytheia; Heracles in the cup was a favorite motif on black-figure pottery.[citation needed] Such a magical conveyance undercuts any literal geography for Erytheia, the "red island" of the sunset.
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+ When Heracles landed at Erytheia, he was confronted by the two-headed dog Orthrus. With one blow from his olive-wood club, Heracles killed Orthrus. Eurytion the herdsman came to assist Orthrus, but Heracles dealt with him the same way.
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+ On hearing the commotion, Geryon sprang into action, carrying three shields and three spears, and wearing three helmets. He attacked Heracles at the River Anthemus, but was slain by one of Heracles' poisoned arrows. Heracles shot so forcefully that the arrow pierced Geryon's forehead, "and Geryon bent his neck over to one side, like a poppy that spoils its delicate shapes, shedding its petals all at once."[24]
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+ Heracles then had to herd the cattle back to Eurystheus. In Roman versions of the narrative, Heracles drove the cattle over the Aventine Hill on the future site of Rome. The giant Cacus, who lived there, stole some of the cattle as Heracles slept, making the cattle walk backwards so that they left no trail, a repetition of the trick of the young Hermes. According to some versions, Heracles drove his remaining cattle past the cave, where Cacus had hidden the stolen animals, and they began calling out to each other. In other versions, Cacus' sister Caca told Heracles where he was. Heracles then killed Cacus, and set up an altar on the spot, later the site of Rome's Forum Boarium (the cattle market).
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+ To annoy Heracles, Hera sent a gadfly to bite the cattle, irritate them, and scatter them. Within a year, Heracles retrieved them. Hera then sent a flood which raised the level of a river so much that Heracles could not cross with the cattle. He piled stones into the river to make the water shallower. When he finally reached the court of Eurystheus, the cattle were sacrificed to Hera.
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+ After Heracles completed the first ten labours, Eurystheus gave him two more, claiming that slaying the Hydra did not count (because Iolaus helped Heracles), neither did cleaning the Augean Stables (either because he was paid for the job or because the rivers did the work).
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+ The first additional labour was to steal three of the golden apples from the garden of the Hesperides. Heracles first caught the Old Man of the Sea, the shapeshifting sea god,[25] to learn where the Garden of the Hesperides was located.[26]
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+ In some variations, Heracles, either at the start or at the end of this task, meets Antaeus, who was invincible as long as he touched his mother, Gaia, the Earth. Heracles killed Antaeus by holding him aloft and crushing him in a bear hug.[27]
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+ Herodotus claims that Heracles stopped in Egypt, where King Busiris decided to make him the yearly sacrifice, but Heracles burst out of his chains.
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+ Heracles finally made his way to the garden of the Hesperides, where he encountered Atlas holding up the heavens on his shoulders. Heracles persuaded Atlas to get the three golden Apples for him by offering to hold up the heavens in his place for a little while. Atlas could get the apples because, in this version, he was the father or otherwise related to the Hesperides. This would have made the labour – like the Hydra and the Augean stables – void because Heracles had received help. When Atlas returned, he decided that he did not want to take the heavens back, and instead offered to deliver the apples himself, but Heracles tricked him by agreeing to remain in place of Atlas on the condition that Atlas relieve him temporarily while Heracles adjusted his cloak. Atlas agreed, but Heracles reneged and walked away with the apples. According to an alternative version, Heracles slew Ladon, the dragon who guarded the apples instead. Eurystheus was furious that Heracles had accomplished something that Eurystheus thought could not possibly be done.
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+ The twelfth and final labour was the capture of Cerberus, the three-headed, dragon-tailed dog that was the guardian of the gates of the Underworld. To prepare for his descent into the Underworld, Heracles went to Eleusis (or Athens) to be initiated in the Eleusinian Mysteries. He entered the Underworld, and Hermes and Athena were his guides.
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+ While in the Underworld, Heracles met Theseus and Pirithous. The two companions had been imprisoned by Hades for attempting to kidnap Persephone. One tradition tells of snakes coiling around their legs, then turning into stone; another that Hades feigned hospitality and prepared a feast inviting them to sit. They unknowingly sat in chairs of forgetfulness and were permanently ensnared. When Heracles had pulled Theseus first from his chair, some of his thigh stuck to it (this explains the supposedly lean thighs of Athenians), but the Earth shook at the attempt to liberate Pirithous, whose desire to have the goddess for himself was so insulting he was doomed to stay behind.
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+ Heracles found Hades and asked permission to bring Cerberus to the surface, which Hades agreed to if Heracles could subdue the beast without using weapons. Heracles overpowered Cerberus with his bare hands and slung the beast over his back. He carried Cerberus out of the Underworld through a cavern entrance in the Peloponnese and brought it to Eurystheus, who again fled into his pithos. Eurystheus begged Heracles to return Cerberus to the Underworld, offering in return to release him from any further labours when Cerberus disappeared back to his master.
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+ After completing the Twelve Labours, one tradition says Heracles joined Jason and the Argonauts in their quest for the Golden Fleece. However, Herodorus (c. 400 BC) disputed this and denied Heracles ever sailed with the Argonauts. A separate tradition (e.g. Argonautica) has Heracles accompany the Argonauts, but he did not travel with them as far as Colchis.
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+ Some ancient Greeks found allegorical meanings of a moral, psychological or philosophical nature in the Labours of Heracles. This trend became more prominent in the Renaissance.[28] For example, Heraclitus the Grammarian wrote in his Homeric Problems:
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+ I turn to Heracles. We must not suppose he attained such power in those days as a result of his physical strength. Rather, he was a man of intellect, an initiate in heavenly wisdom, who, as it were, shed light on philosophy, which had been hidden in deep darkness. The most authoritative of the Stoics agree with this account.... The (Erymanthian) boar which he overcame is the common incontinence of men; the (Nemean) lion is the indiscriminate rush towards improper goals; in the same way, by fettering irrational passions he gave rise to the belief that he had fettered the violent (Cretan) bull. He banished cowardice also from the world, in the shape of the hind of Ceryneia. There was another "labor" too, not properly so called, in which he cleared out the mass of dung (from the Augean stables) — in other words, the foulness that disfigures humanity. The (Stymphalian) birds he scattered are the windy hopes that feed our lives; the many-headed hydra that he burned, as it were, with the fires of exhortation, is pleasure, which begins to grow again as soon as it is cut out.
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1
+ The Twelve Labours of Heracles or Hercules (Greek: οἱ Ἡρακλέους ἆθλοι, hoi Hērakléous âthloi)[1][2] are a series of episodes concerning a penance carried out by Heracles, the greatest of the Greek heroes, whose name was later romanised as Hercules. They were accomplished at the service of King Eurystheus. The episodes were later connected by a continuous narrative. The establishment of a fixed cycle of twelve labours was attributed by the Greeks to an epic poem, now lost, written by Peisander, dated about 600 BC.[3] After Heracles killed his wife and children, he went to the oracle at Delphi. He prayed to the god Apollo for guidance. Heracles was told to serve the king of Mycenae, Eurystheus, for ten years. During this time, he is sent to perform a series of difficult feats, called labours.[4]
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+ Driven mad by Hera (queen of the gods), Heracles slew his sons by his wife Megara.[5] After recovering his sanity, Heracles deeply regretted his actions; he was purified by King Thespius, then traveled to Delphi to inquire how he could atone for his actions. Pythia, the Oracle of Delphi, advised him to go to Tiryns and serve his cousin, King Eurystheus, for ten years, performing whatever labours Eurystheus might set him; in return, he would be rewarded with immortality. Heracles despaired at this, loathing to serve a man whom he knew to be far inferior to himself, yet fearing to oppose his father Zeus. Eventually, he placed himself at Eurystheus's disposal.
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+ Eurystheus originally ordered Heracles to perform ten labours. Heracles accomplished these tasks, but Eurystheus refused to recognize two: the slaying of the Lernaean Hydra, as Heracles' nephew and charioteer Iolaus had helped him; and the cleansing of the Augeas, because Heracles accepted payment for the labour. Eurystheus set two more tasks (fetching the Golden Apples of Hesperides and capturing Cerberus), which Heracles also performed, bringing the total number of tasks to twelve.
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+ As they survive, the labours of Heracles are not recounted in any single place, but must be reassembled from many sources. Ruck and Staples[6] assert that there is no one way to interpret the labours, but that six were located in the Peloponnese, culminating with the rededication of Olympia. Six others took the hero farther afield, to places that were, per Ruck, "all previously strongholds of Hera or the 'Goddess' and were Entrances to the Netherworld".[6] In each case, the pattern was the same: Heracles was sent to kill or subdue, or to fetch back for Eurystheus (as Hera's representative) a magical animal or plant.
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+ A famous depiction of the labours in Greek sculpture is found on the metopes of the Temple of Zeus at Olympia, which date to the 450s BC.[citation needed]
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+ In his labours, Heracles was sometimes accompanied by a male companion (an eromenos), according to Licymnius[citation needed] and others, such as Iolaus, his nephew. Although he was supposed to perform only ten labours, this assistance led to two labours being disqualified: Eurystheus refused to recognize slaying the Hydra, because Iolaus helped him, and the cleansing of the Augean stables, because Heracles was paid for his services and/or because the rivers did the work. Several of the labours involved the offspring (by various accounts) of Typhon and his mate Echidna, all overcome by Heracles.
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+ A traditional order of the labours found in the Bibliotheca[7] by Pseudo-Apollodorus is:
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+ Heracles wandered the area until he came to the town of Cleonae. There he met a boy who said that if Heracles slew the Nemean lion and returned alive within 30 days, the town would sacrifice a lion to Zeus, but if he did not return within 30 days or if he died, the boy would sacrifice himself to Zeus. Another version claims that he met Molorchos, a shepherd who had lost his son to the lion, saying that if he came back within 30 days, a ram would be sacrificed to Zeus. If he did not return within 30 days, it would be sacrificed to the dead Heracles as a mourning offering.
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+ While searching for the lion, Heracles fletched some arrows to use against it, not knowing that its golden fur was impenetrable. When he found and shot the lion, firing at it with his bow, he discovered the fur's protective property as the arrow bounced harmlessly off the creature's thigh. After some time, Heracles made the lion return to his cave. The cave had two entrances, one of which Heracles blocked; he then entered the other. In those dark and close quarters, Heracles stunned the beast with his club and, using his immense strength, strangled it to death. During the fight the lion bit off one of his fingers.[8] Others say that he shot arrows at it, eventually shooting it in the unarmored mouth. After slaying the lion, he tried to skin it with a knife from his belt, but failed. He then tried sharpening the knife with a stone and even tried with the stone itself. Finally, Athena, noticing the hero's plight, told Heracles to use one of the lion's own claws to skin the pelt. Others say that Heracles' armor was, in fact, the hide of the Lion of Cithaeron.
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+ When he returned on the 30th day carrying the carcass of the lion on his shoulders, King Eurystheus was amazed and terrified. Eurystheus forbade him ever again to enter the city; from then on he was to display the fruits of his labours outside the city gates. Eurystheus would then tell Heracles his tasks through a herald, not personally. Eurystheus even had a large bronze jar made for him in which to hide from Heracles if need be. Eurystheus then warned him that the tasks would become increasingly difficult.
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+ Heracles' second labour was to slay the Lernaean Hydra, which Hera had raised just to slay Heracles. Upon reaching the swamp near Lake Lerna, where the Hydra dwelt, Heracles used a cloth to cover his mouth and nose to protect himself from the poisonous fumes. He fired flaming arrows into the Hydra's lair, the spring of Amymone, a deep cave that it only came out of to terrorize neighboring villages.[9] He then confronted the Hydra, wielding a harvesting sickle (according to some early vase-paintings), a sword or his famed club. Ruck and Staples (1994: 170) have pointed out that the chthonic creature's reaction was botanical: upon cutting off each of its heads he found that two grew back, an expression of the hopelessness of such a struggle for any but the hero. Additionally, one of the Hydra's heads - the middle one - was immortal.
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+ The details of the struggle are explicit in the Bibliotheca (2.5.2): realizing that he could not defeat the Hydra in this way, Heracles called on his nephew Iolaus for help. His nephew then came upon the idea (possibly inspired by Athena) of using a firebrand to scorch the neck stumps after each decapitation. Heracles cut off each head and Iolaus cauterized the open stumps. Seeing that Heracles was winning the struggle, Hera sent a giant crab to distract him. He crushed it under his mighty foot. He cut off the Hydra's one immortal head with a golden sword given to him by Athena. Heracles placed it under a great rock on the sacred way between Lerna and Elaius (Kerenyi 1959:144), and dipped his arrows in the Hydra's poisonous blood, and so his second task was complete. The alternative version of this myth is that after cutting off one head, he then dipped his sword in it and used its venom to burn each head so it could not grow back. Hera, upset that Heracles had slain the beast she raised to kill him, placed it in the dark blue vault of the sky as the constellation Hydra. She then turned the crab into the constellation Cancer.
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+ Later, Heracles used an arrow dipped in the Hydra's poisonous blood to kill the centaur Nessus; and Nessus's tainted blood was applied to the Tunic of Nessus, by which the centaur had his posthumous revenge. Both Strabo and Pausanias report that the stench of the river Anigrus in Elis, making all the fish of the river inedible, was reputed to be due to the Hydra's venom, washed from the arrows Heracles used on the centaur.[10]
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+ Eurystheus and Hera were greatly angered that Heracles had survived the Nemean Lion and the Lernaean Hydra. For the third labour, they found a task which they thought would spell doom for the hero. It was not slaying a beast or monster, as it had already been established that Heracles could overcome even the most fearsome opponents. Instead, Eurystheus ordered him to capture the Ceryneian Hind, which was so fast that it could outrun an arrow.
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+ After beginning the search, Heracles awoke from sleeping and saw the hind by the glint on its antlers. Heracles then chased the hind on foot for a full year through Greece, Thrace, Istria, and the land of the Hyperboreans. In some versions, he captured the hind while it slept, rendering it lame with a trap net. In other versions, he encountered Artemis in her temple; she told him to leave the hind and tell Eurystheus all that had happened, and his third labour would be considered to be completed. Yet another version claims that Heracles trapped the Hind with an arrow between its forelegs.
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+ Eurystheus had given Heracles this task hoping to incite Artemis' anger at Heracles for his desecration of her sacred animal. As he was returning with the hind, Heracles encountered Artemis and her brother Apollo. He begged the goddess for forgiveness, explaining that he had to catch it as part of his penance, but he promised to return it. Artemis forgave him, foiling Eurystheus' plan to have her punish him.
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+ Upon bringing the hind to Eurystheus, he was told that it was to become part of the King's menagerie. Heracles knew that he had to return the hind as he had promised, so he agreed to hand it over on the condition that Eurystheus himself come out and take it from him. The King came out, but the moment that Heracles let the hind go, it sprinted back to its mistress and Heracles left, saying that Eurystheus had not been quick enough.
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+ Eurystheus was disappointed that Heracles had overcome yet another creature and was humiliated by the hind's escape, so he assigned Heracles another dangerous task. By some accounts, the fourth labour was to bring the fearsome Erymanthian Boar back to Eurystheus alive (there is no single definitive telling of the labours). On the way to Mount Erymanthos where the boar lived, Heracles visited Pholus ("caveman"), a kind and hospitable centaur and old friend. Heracles ate with Pholus in his cavern (though the centaur devoured his meat raw) and asked for wine. Pholus had only one jar of wine, a gift from Dionysus to all the centaurs on Mount Erymanthos. Heracles convinced him to open it, and the smell attracted the other centaurs. They did not understand that wine needs to be tempered with water, became drunk, and attacked Heracles. Heracles shot at them with his poisonous arrows, killing many, and the centaurs retreated all the way to Chiron's cave.
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+ Pholus was curious why the arrows caused so much death. He picked one up but dropped it, and the arrow stabbed his hoof, poisoning him. One version states that a stray arrow hit Chiron as well. He was immortal, but he still felt the pain. Chiron's pain was so great that he volunteered to give up his immortality and take the place of Prometheus, who had been chained to the top of a mountain to have his liver eaten daily by an eagle. Prometheus' torturer, the eagle, continued its torture on Chiron, so Heracles shot it dead with an arrow. It is generally accepted that the tale was meant to show Heracles as being the recipient of Chiron's surrendered immortality. However, this tale contradicts the fact that Chiron later taught Achilles. The tale of the centaurs sometimes appears in other parts of the twelve labours, as does the freeing of Prometheus.
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+ Heracles had visited Chiron to gain advice on how to catch the boar, and Chiron had told him to drive it into thick snow, which sets this labour in mid-winter. Heracles caught the boar, bound it, and carried it back to Eurystheus, who was frightened of it and ducked down in his half-buried storage pithos, begging Heracles to get rid of the beast.
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+ The fifth labour was to clean the stables of King Augeas. This assignment was intended to be both humiliating (rather than impressive, as the previous labours had been) and impossible, since the livestock were divinely healthy (and immortal) and therefore produced an enormous quantity of dung. The Augean Stables (/ɔːˈdʒiːən/) had not been cleaned in over 30 years, and over 1,000 cattle lived there. However, Heracles succeeded by re-routing the rivers Alpheus and Peneus to wash out the filth.
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+ Before starting on the task, Heracles had asked Augeas for one-tenth of the cattle if he finished the task in one day, and Augeas agreed. But afterwards Augeas refused to honour the agreement on the grounds that Heracles had been ordered to carry out the task by Eurystheus anyway. Heracles claimed his reward in court, and was supported by Augeas' son Phyleus. Augeas banished them both before the court had ruled. Heracles returned, slew Augeas, and gave his kingdom to Phyleus. Heracles then founded the Olympic Games.
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+ The success of this labour was ultimately discounted as the rushing waters had done the work of cleaning the stables and because Heracles was paid for doing the labour.
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+ Eurystheus said that Heracles still had seven labours to perform.[11]
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+ The sixth labour was to defeat the Stymphalian birds, man-eating birds with beaks made of bronze and sharp metallic feathers they could launch at their victims. They were sacred to Ares, the god of war. Furthermore, their dung was highly toxic. They had migrated to Lake Stymphalia in Arcadia, where they bred quickly and took over the countryside, destroying local crops, fruit trees, and townspeople. Heracles could not go too far into the swamp, for it would not support his weight. Athena, noticing the hero's plight, gave Heracles a rattle which Hephaestus had made especially for the occasion. Heracles shook the rattle and frightened the birds into the air. Heracles then shot many of them with his arrows. The rest flew far away, never to return. The Argonauts would later encounter them.
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+ The seventh labour was to capture the Cretan Bull, father of the Minotaur. Heracles sailed to Crete, where King Minos gave Heracles permission to take the bull away and even offered him assistance (which Heracles declined, plausibly because he did not want the labour to be discounted as before).[12] The bull had been wreaking havoc on Crete by uprooting crops and leveling orchard walls. Heracles sneaked up behind the bull and then used his hands to throttle it (stopping before it was killed), and then shipped it back to Tiryns. Eurystheus, who hid in his pithos at first sight of the creature, wanted to sacrifice the bull to Hera, who hated Heracles. She refused the sacrifice because it reflected glory on Heracles. The bull was released and wandered into Marathon, becoming known as the Marathonian Bull.[12] Theseus would later sacrifice the bull to Athena and/or Apollo.
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+ As the eighth of his Twelve Labours, also categorised as the second of the Non-Peloponneisan labours,[13] Heracles was sent by King Eurystheus to steal the Mares from Diomedes. The mares’ madness was attributed to their unnatural diet which consisted of the flesh[14] of unsuspecting guests or strangers to the island.[15] Some versions of the myth say that the mares also expelled fire when they breathed.[16] The Mares, which were the terror of Thrace, were kept tethered by iron chains to a bronze manger in the now vanished city of Tirida[17] and were named Podargos (the swift), Lampon (the shining), Xanthos (the yellow) and Deinos (or Deinus, the terrible).[18] Although very similar, there are slight variances in the exact details regarding the mares’ capture.
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+ In one version, Heracles brought a number of volunteers to help him capture the giant horses.[17] After overpowering Diomedes’ men, Heracles broke the chains that tethered the horses and drove the mares down to sea. Unaware that the mares were man-eating and uncontrollable, Heracles left them in the charge of his favored companion, Abderus, while he left to fight Diomedes. Upon his return, Heracles found that the boy was eaten. As revenge, Heracles fed Diomedes to his own horses and then founded Abdera next to the boy's tomb.[15]
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+ In another version, Heracles, who was visiting the island, stayed awake so that he didn't have his throat cut by Diomedes in the night, and cut the chains binding the horses once everyone was asleep. Having scared the horses onto the high ground of a knoll, Heracles quickly dug a trench through the peninsula, filling it with water and thus flooding the low-lying plain. When Diomedes and his men turned to flee, Heracles killed them with an axe (or a club[17]), and fed Diomedes’ body to the horses to calm them.
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+ In yet another version, Heracles first captured Diomedes and fed him to the mares before releasing them. Only after realizing that their King was dead did his men, the Bistonians,[15][17] attack Heracles. Upon seeing the mares charging at them, led in a chariot by Abderus, the Bistonians turned and fled.
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+ All versions have eating human flesh make the horses calmer, giving Heracles the opportunity to bind their mouths shut, and easily take them back to King Eurystheus, who dedicated the horses to Hera.[19] In some versions, they were allowed to roam freely around Argos, having become permanently calm, but in others, Eurystheus ordered the horses taken to Olympus to be sacrificed to Zeus, but Zeus refused them, and sent wolves, lions, and bears to kill them.[20] Roger Lancelyn Green states in his Tales of the Greek Heroes that the mares’ descendants were used in the Trojan War, and survived even to the time of Alexander the Great.[17][21] After the incident, Eurystheus sent Heracles to bring back Hippolyta's Girdle.
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+ Eurystheus' daughter Admete wanted the Belt of Hippolyta, queen of the Amazons, a gift from her father Ares. To please his daughter, Eurystheus ordered Heracles to retrieve the belt as his ninth labour.
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+ Taking a band of friends with him, Heracles set sail, stopping at the island of Paros, which was inhabited by some sons of Minos. The sons killed two of Heracles' companions, an act which set Heracles on a rampage. He killed two of the sons of Minos and threatened the other inhabitants until he was offered two men to replace his fallen companions. Heracles agreed and took two of Minos' grandsons, Alcaeus and Sthenelus. They continued their voyage and landed at the court of Lycus, whom Heracles defended in a battle against King Mygdon of Bebryces. After killing King Mygdon, Heracles gave much of the land to his friend Lycus. Lycus called the land Heraclea. The crew then set off for Themiscyra, where Hippolyta lived.
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+ All would have gone well for Heracles had it not been for Hera. Hippolyta, impressed with Heracles and his exploits, agreed to give him the belt and would have done so had Hera not disguised herself and walked among the Amazons sowing seeds of distrust. She claimed the strangers were plotting to carry off the queen of the Amazons. Alarmed, the women set off on horseback to confront Heracles. When Heracles saw them, he thought Hippolyta had been plotting such treachery all along and had never meant to hand over the belt, so he killed her, took the belt and returned to Eurystheus.
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+ The tenth labour was to obtain the Cattle of the three-bodied giant Geryon. In the fullest account in the Bibliotheca of Pseudo-Apollodorus,[22] Heracles had to go to the island of Erytheia in the far west (sometimes identified with the Hesperides, or with the island which forms the city of Cádiz) to get the cattle. On the way there, he crossed the Libyan desert[23] and became so frustrated at the heat that he shot an arrow at the Sun. The sun-god Helios "in admiration of his courage" gave Heracles the golden cup Helios used to sail across the sea from west to east each night. Heracles rode the cup to Erytheia; Heracles in the cup was a favorite motif on black-figure pottery.[citation needed] Such a magical conveyance undercuts any literal geography for Erytheia, the "red island" of the sunset.
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+ When Heracles landed at Erytheia, he was confronted by the two-headed dog Orthrus. With one blow from his olive-wood club, Heracles killed Orthrus. Eurytion the herdsman came to assist Orthrus, but Heracles dealt with him the same way.
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+ On hearing the commotion, Geryon sprang into action, carrying three shields and three spears, and wearing three helmets. He attacked Heracles at the River Anthemus, but was slain by one of Heracles' poisoned arrows. Heracles shot so forcefully that the arrow pierced Geryon's forehead, "and Geryon bent his neck over to one side, like a poppy that spoils its delicate shapes, shedding its petals all at once."[24]
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+ Heracles then had to herd the cattle back to Eurystheus. In Roman versions of the narrative, Heracles drove the cattle over the Aventine Hill on the future site of Rome. The giant Cacus, who lived there, stole some of the cattle as Heracles slept, making the cattle walk backwards so that they left no trail, a repetition of the trick of the young Hermes. According to some versions, Heracles drove his remaining cattle past the cave, where Cacus had hidden the stolen animals, and they began calling out to each other. In other versions, Cacus' sister Caca told Heracles where he was. Heracles then killed Cacus, and set up an altar on the spot, later the site of Rome's Forum Boarium (the cattle market).
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+ To annoy Heracles, Hera sent a gadfly to bite the cattle, irritate them, and scatter them. Within a year, Heracles retrieved them. Hera then sent a flood which raised the level of a river so much that Heracles could not cross with the cattle. He piled stones into the river to make the water shallower. When he finally reached the court of Eurystheus, the cattle were sacrificed to Hera.
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+ After Heracles completed the first ten labours, Eurystheus gave him two more, claiming that slaying the Hydra did not count (because Iolaus helped Heracles), neither did cleaning the Augean Stables (either because he was paid for the job or because the rivers did the work).
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+ The first additional labour was to steal three of the golden apples from the garden of the Hesperides. Heracles first caught the Old Man of the Sea, the shapeshifting sea god,[25] to learn where the Garden of the Hesperides was located.[26]
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+ In some variations, Heracles, either at the start or at the end of this task, meets Antaeus, who was invincible as long as he touched his mother, Gaia, the Earth. Heracles killed Antaeus by holding him aloft and crushing him in a bear hug.[27]
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+ Herodotus claims that Heracles stopped in Egypt, where King Busiris decided to make him the yearly sacrifice, but Heracles burst out of his chains.
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+ Heracles finally made his way to the garden of the Hesperides, where he encountered Atlas holding up the heavens on his shoulders. Heracles persuaded Atlas to get the three golden Apples for him by offering to hold up the heavens in his place for a little while. Atlas could get the apples because, in this version, he was the father or otherwise related to the Hesperides. This would have made the labour – like the Hydra and the Augean stables – void because Heracles had received help. When Atlas returned, he decided that he did not want to take the heavens back, and instead offered to deliver the apples himself, but Heracles tricked him by agreeing to remain in place of Atlas on the condition that Atlas relieve him temporarily while Heracles adjusted his cloak. Atlas agreed, but Heracles reneged and walked away with the apples. According to an alternative version, Heracles slew Ladon, the dragon who guarded the apples instead. Eurystheus was furious that Heracles had accomplished something that Eurystheus thought could not possibly be done.
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+ The twelfth and final labour was the capture of Cerberus, the three-headed, dragon-tailed dog that was the guardian of the gates of the Underworld. To prepare for his descent into the Underworld, Heracles went to Eleusis (or Athens) to be initiated in the Eleusinian Mysteries. He entered the Underworld, and Hermes and Athena were his guides.
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+ While in the Underworld, Heracles met Theseus and Pirithous. The two companions had been imprisoned by Hades for attempting to kidnap Persephone. One tradition tells of snakes coiling around their legs, then turning into stone; another that Hades feigned hospitality and prepared a feast inviting them to sit. They unknowingly sat in chairs of forgetfulness and were permanently ensnared. When Heracles had pulled Theseus first from his chair, some of his thigh stuck to it (this explains the supposedly lean thighs of Athenians), but the Earth shook at the attempt to liberate Pirithous, whose desire to have the goddess for himself was so insulting he was doomed to stay behind.
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+ Heracles found Hades and asked permission to bring Cerberus to the surface, which Hades agreed to if Heracles could subdue the beast without using weapons. Heracles overpowered Cerberus with his bare hands and slung the beast over his back. He carried Cerberus out of the Underworld through a cavern entrance in the Peloponnese and brought it to Eurystheus, who again fled into his pithos. Eurystheus begged Heracles to return Cerberus to the Underworld, offering in return to release him from any further labours when Cerberus disappeared back to his master.
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+ After completing the Twelve Labours, one tradition says Heracles joined Jason and the Argonauts in their quest for the Golden Fleece. However, Herodorus (c. 400 BC) disputed this and denied Heracles ever sailed with the Argonauts. A separate tradition (e.g. Argonautica) has Heracles accompany the Argonauts, but he did not travel with them as far as Colchis.
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+ Some ancient Greeks found allegorical meanings of a moral, psychological or philosophical nature in the Labours of Heracles. This trend became more prominent in the Renaissance.[28] For example, Heraclitus the Grammarian wrote in his Homeric Problems:
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+ I turn to Heracles. We must not suppose he attained such power in those days as a result of his physical strength. Rather, he was a man of intellect, an initiate in heavenly wisdom, who, as it were, shed light on philosophy, which had been hidden in deep darkness. The most authoritative of the Stoics agree with this account.... The (Erymanthian) boar which he overcame is the common incontinence of men; the (Nemean) lion is the indiscriminate rush towards improper goals; in the same way, by fettering irrational passions he gave rise to the belief that he had fettered the violent (Cretan) bull. He banished cowardice also from the world, in the shape of the hind of Ceryneia. There was another "labor" too, not properly so called, in which he cleared out the mass of dung (from the Augean stables) — in other words, the foulness that disfigures humanity. The (Stymphalian) birds he scattered are the windy hopes that feed our lives; the many-headed hydra that he burned, as it were, with the fires of exhortation, is pleasure, which begins to grow again as soon as it is cut out.
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+ An earthquake (also known as a quake, tremor or temblor) is the shaking of the surface of the Earth resulting from a sudden release of energy in the Earth's lithosphere that creates seismic waves. Earthquakes can range in size from those that are so weak that they cannot be felt to those violent enough to propel objects and people into the air, and wreak destruction across entire cities. The seismicity, or seismic activity, of an area is the frequency, type, and size of earthquakes experienced over a period of time. The word tremor is also used for non-earthquake seismic rumbling.
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+ At the Earth's surface, earthquakes manifest themselves by shaking and displacing or disrupting the ground. When the epicenter of a large earthquake is located offshore, the seabed may be displaced sufficiently to cause a tsunami. Earthquakes can also trigger landslides and occasionally, volcanic activity.
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+ In its most general sense, the word earthquake is used to describe any seismic event—whether natural or caused by humans—that generates seismic waves. Earthquakes are caused mostly by rupture of geological faults but also by other events such as volcanic activity, landslides, mine blasts, and nuclear tests. An earthquake's point of initial rupture is called its hypocenter or focus. The epicenter is the point at ground level directly above the hypocenter.
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+ Tectonic earthquakes occur anywhere in the earth where there is sufficient stored elastic strain energy to drive fracture propagation along a fault plane. The sides of a fault move past each other smoothly and aseismically only if there are no irregularities or asperities along the fault surface that increase the frictional resistance. Most fault surfaces do have such asperities, which leads to a form of stick-slip behavior. Once the fault has locked, continued relative motion between the plates leads to increasing stress and therefore, stored strain energy in the volume around the fault surface. This continues until the stress has risen sufficiently to break through the asperity, suddenly allowing sliding over the locked portion of the fault, releasing the stored energy.[1] This energy is released as a combination of radiated elastic strain seismic waves,[2] frictional heating of the fault surface, and cracking of the rock, thus causing an earthquake. This process of gradual build-up of strain and stress punctuated by occasional sudden earthquake failure is referred to as the elastic-rebound theory. It is estimated that only 10 percent or less of an earthquake's total energy is radiated as seismic energy. Most of the earthquake's energy is used to power the earthquake fracture growth or is converted into heat generated by friction. Therefore, earthquakes lower the Earth's available elastic potential energy and raise its temperature, though these changes are negligible compared to the conductive and convective flow of heat out from the Earth's deep interior.[3]
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+ There are three main types of fault, all of which may cause an interplate earthquake: normal, reverse (thrust), and strike-slip. Normal and reverse faulting are examples of dip-slip, where the displacement along the fault is in the direction of dip and where movement on them involves a vertical component. Normal faults occur mainly in areas where the crust is being extended such as a divergent boundary. Reverse faults occur in areas where the crust is being shortened such as at a convergent boundary. Strike-slip faults are steep structures where the two sides of the fault slip horizontally past each other; transform boundaries are a particular type of strike-slip fault. Many earthquakes are caused by movement on faults that have components of both dip-slip and strike-slip; this is known as oblique slip.
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+ Reverse faults, particularly those along convergent plate boundaries, are associated with the most powerful earthquakes, megathrust earthquakes, including almost all of those of magnitude 8 or more. Strike-slip faults, particularly continental transforms, can produce major earthquakes up to about magnitude 8. Earthquakes associated with normal faults are generally less than magnitude 7. For every unit increase in magnitude, there is a roughly thirtyfold increase in the energy released. For instance, an earthquake of magnitude 6.0 releases approximately 32 times more energy than a 5.0 magnitude earthquake and a 7.0 magnitude earthquake releases 1,000 times more energy than a 5.0 magnitude of earthquake. An 8.6 magnitude earthquake releases the same amount of energy as 10,000 atomic bombs like those used in World War II.[4]
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+ This is so because the energy released in an earthquake, and thus its magnitude, is proportional to the area of the fault that ruptures[5] and the stress drop. Therefore, the longer the length and the wider the width of the faulted area, the larger the resulting magnitude. The topmost, brittle part of the Earth's crust, and the cool slabs of the tectonic plates that are descending down into the hot mantle, are the only parts of our planet that can store elastic energy and release it in fault ruptures. Rocks hotter than about 300 °C (572 °F) flow in response to stress; they do not rupture in earthquakes.[6][7] The maximum observed lengths of ruptures and mapped faults (which may break in a single rupture) are approximately 1,000 km (620 mi). Examples are the earthquakes in Alaska (1957), Chile (1960), and Sumatra (2004), all in subduction zones. The longest earthquake ruptures on strike-slip faults, like the San Andreas Fault (1857, 1906), the North Anatolian Fault in Turkey (1939), and the Denali Fault in Alaska (2002), are about half to one third as long as the lengths along subducting plate margins, and those along normal faults are even shorter.
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+ The most important parameter controlling the maximum earthquake magnitude on a fault, however, is not the maximum available length, but the available width because the latter varies by a factor of 20. Along converging plate margins, the dip angle of the rupture plane is very shallow, typically about 10 degrees.[8] Thus, the width of the plane within the top brittle crust of the Earth can become 50–100 km (31–62 mi) (Japan, 2011; Alaska, 1964), making the most powerful earthquakes possible.
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+ Strike-slip faults tend to be oriented near vertically, resulting in an approximate width of 10 km (6.2 mi) within the brittle crust.[9] Thus, earthquakes with magnitudes much larger than 8 are not possible. Maximum magnitudes along many normal faults are even more limited because many of them are located along spreading centers, as in Iceland, where the thickness of the brittle layer is only about six kilometres (3.7 mi).[10][11]
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+ In addition, there exists a hierarchy of stress level in the three fault types. Thrust faults are generated by the highest, strike-slip by intermediate, and normal faults by the lowest stress levels.[12] This can easily be understood by considering the direction of the greatest principal stress, the direction of the force that "pushes" the rock mass during the faulting. In the case of normal faults, the rock mass is pushed down in a vertical direction, thus the pushing force (greatest principal stress) equals the weight of the rock mass itself. In the case of thrusting, the rock mass "escapes" in the direction of the least principal stress, namely upward, lifting the rock mass up, and thus, the overburden equals the least principal stress. Strike-slip faulting is intermediate between the other two types described above. This difference in stress regime in the three faulting environments can contribute to differences in stress drop during faulting, which contributes to differences in the radiated energy, regardless of fault dimensions.
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+ Where plate boundaries occur within the continental lithosphere, deformation is spread out over a much larger area than the plate boundary itself. In the case of the San Andreas fault continental transform, many earthquakes occur away from the plate boundary and are related to strains developed within the broader zone of deformation caused by major irregularities in the fault trace (e.g., the "Big bend" region). The Northridge earthquake was associated with movement on a blind thrust within such a zone. Another example is the strongly oblique convergent plate boundary between the Arabian and Eurasian plates where it runs through the northwestern part of the Zagros Mountains. The deformation associated with this plate boundary is partitioned into nearly pure thrust sense movements perpendicular to the boundary over a wide zone to the southwest and nearly pure strike-slip motion along the Main Recent Fault close to the actual plate boundary itself. This is demonstrated by earthquake focal mechanisms.[13]
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+ All tectonic plates have internal stress fields caused by their interactions with neighboring plates and sedimentary loading or unloading (e.g., deglaciation).[14] These stresses may be sufficient to cause failure along existing fault planes, giving rise to intraplate earthquakes.[15]
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+ The majority of tectonic earthquakes originate at the ring of fire in depths not exceeding tens of kilometers. Earthquakes occurring at a depth of less than 70 km (43 mi) are classified as "shallow-focus" earthquakes, while those with a focal-depth between 70 and 300 km (43 and 186 mi) are commonly termed "mid-focus" or "intermediate-depth" earthquakes. In subduction zones, where older and colder oceanic crust descends beneath another tectonic plate, deep-focus earthquakes may occur at much greater depths (ranging from 300 to 700 km (190 to 430 mi)).[16] These seismically active areas of subduction are known as Wadati–Benioff zones. Deep-focus earthquakes occur at a depth where the subducted lithosphere should no longer be brittle, due to the high temperature and pressure. A possible mechanism for the generation of deep-focus earthquakes is faulting caused by olivine undergoing a phase transition into a spinel structure.[17]
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+ Earthquakes often occur in volcanic regions and are caused there, both by tectonic faults and the movement of magma in volcanoes. Such earthquakes can serve as an early warning of volcanic eruptions, as during the 1980 eruption of Mount St. Helens.[18] Earthquake swarms can serve as markers for the location of the flowing magma throughout the volcanoes. These swarms can be recorded by seismometers and tiltmeters (a device that measures ground slope) and used as sensors to predict imminent or upcoming eruptions.[19]
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+ A tectonic earthquake begins by an initial rupture at a point on the fault surface, a process known as nucleation. The scale of the nucleation zone is uncertain, with some evidence, such as the rupture dimensions of the smallest earthquakes, suggesting that it is smaller than 100 m (330 ft) while other evidence, such as a slow component revealed by low-frequency spectra of some earthquakes, suggest that it is larger. The possibility that the nucleation involves some sort of preparation process is supported by the observation that about 40% of earthquakes are preceded by foreshocks. Once the rupture has initiated, it begins to propagate along the fault surface. The mechanics of this process are poorly understood, partly because it is difficult to recreate the high sliding velocities in a laboratory. Also the effects of strong ground motion make it very difficult to record information close to a nucleation zone.[20]
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+ Rupture propagation is generally modeled using a fracture mechanics approach, likening the rupture to a propagating mixed mode shear crack. The rupture velocity is a function of the fracture energy in the volume around the crack tip, increasing with decreasing fracture energy. The velocity of rupture propagation is orders of magnitude faster than the displacement velocity across the fault. Earthquake ruptures typically propagate at velocities that are in the range 70–90% of the S-wave velocity, which is independent of earthquake size. A small subset of earthquake ruptures appear to have propagated at speeds greater than the S-wave velocity. These supershear earthquakes have all been observed during large strike-slip events. The unusually wide zone of coseismic damage caused by the 2001 Kunlun earthquake has been attributed to the effects of the sonic boom developed in such earthquakes. Some earthquake ruptures travel at unusually low velocities and are referred to as slow earthquakes. A particularly dangerous form of slow earthquake is the tsunami earthquake, observed where the relatively low felt intensities, caused by the slow propagation speed of some great earthquakes, fail to alert the population of the neighboring coast, as in the 1896 Sanriku earthquake.[20]
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+ Tides may induce some seismicity. See tidal triggering of earthquakes for details.
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+ Most earthquakes form part of a sequence, related to each other in terms of location and time.[21] Most earthquake clusters consist of small tremors that cause little to no damage, but there is a theory that earthquakes can recur in a regular pattern.[22]
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+ An aftershock is an earthquake that occurs after a previous earthquake, the mainshock. An aftershock is in the same region of the main shock but always of a smaller magnitude. If an aftershock is larger than the main shock, the aftershock is redesignated as the main shock and the original main shock is redesignated as a foreshock. Aftershocks are formed as the crust around the displaced fault plane adjusts to the effects of the main shock.[21]
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+ Earthquake swarms are sequences of earthquakes striking in a specific area within a short period of time. They are different from earthquakes followed by a series of aftershocks by the fact that no single earthquake in the sequence is obviously the main shock, so none has a notable higher magnitude than another. An example of an earthquake swarm is the 2004 activity at Yellowstone National Park.[23] In August 2012, a swarm of earthquakes shook Southern California's Imperial Valley, showing the most recorded activity in the area since the 1970s.[24]
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+ Sometimes a series of earthquakes occur in what has been called an earthquake storm, where the earthquakes strike a fault in clusters, each triggered by the shaking or stress redistribution of the previous earthquakes. Similar to aftershocks but on adjacent segments of fault, these storms occur over the course of years, and with some of the later earthquakes as damaging as the early ones. Such a pattern was observed in the sequence of about a dozen earthquakes that struck the North Anatolian Fault in Turkey in the 20th century and has been inferred for older anomalous clusters of large earthquakes in the Middle East.[25][26]
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+ Quaking or shaking of the earth is a common phenomenon undoubtedly known to humans from earliest times. Prior to the development of strong-motion accelerometers that can measure peak ground speed and acceleration directly, the intensity of the earth-shaking was estimated on the basis of the observed effects, as categorized on various seismic intensity scales. Only in the last century has the source of such shaking been identified as ruptures in the Earth's crust, with the intensity of shaking at any locality dependent not only on the local ground conditions but also on the strength or magnitude of the rupture, and on its distance.[27]
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+ The first scale for measuring earthquake magnitudes was developed by Charles F. Richter in 1935. Subsequent scales (see seismic magnitude scales) have retained a key feature, where each unit represents a ten-fold difference in the amplitude of the ground shaking and a 32-fold difference in energy. Subsequent scales are also adjusted to have approximately the same numeric value within the limits of the scale.[28]
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+ Although the mass media commonly reports earthquake magnitudes as "Richter magnitude" or "Richter scale", standard practice by most seismological authorities is to express an earthquake's strength on the moment magnitude scale, which is based on the actual energy released by an earthquake.[29]
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+ It is estimated that around 500,000 earthquakes occur each year, detectable with current instrumentation. About 100,000 of these can be felt.[30][31] Minor earthquakes occur nearly constantly around the world in places like California and Alaska in the U.S., as well as in El Salvador, Mexico, Guatemala, Chile, Peru, Indonesia, Philippines, Iran, Pakistan, the Azores in Portugal, Turkey, New Zealand, Greece, Italy, India, Nepal and Japan.[32] Larger earthquakes occur less frequently, the relationship being exponential; for example, roughly ten times as many earthquakes larger than magnitude 4 occur in a particular time period than earthquakes larger than magnitude 5.[33] In the (low seismicity) United Kingdom, for example, it has been calculated that the average recurrences are:
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+ an earthquake of 3.7–4.6 every year, an earthquake of 4.7–5.5 every 10 years, and an earthquake of 5.6 or larger every 100 years.[34] This is an example of the Gutenberg–Richter law.
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+ The number of seismic stations has increased from about 350 in 1931 to many thousands today. As a result, many more earthquakes are reported than in the past, but this is because of the vast improvement in instrumentation, rather than an increase in the number of earthquakes. The United States Geological Survey estimates that, since 1900, there have been an average of 18 major earthquakes (magnitude 7.0–7.9) and one great earthquake (magnitude 8.0 or greater) per year, and that this average has been relatively stable.[36] In recent years, the number of major earthquakes per year has decreased, though this is probably a statistical fluctuation rather than a systematic trend.[37] More detailed statistics on the size and frequency of earthquakes is available from the United States Geological Survey (USGS).[38]
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+ A recent increase in the number of major earthquakes has been noted, which could be explained by a cyclical pattern of periods of intense tectonic activity, interspersed with longer periods of low intensity. However, accurate recordings of earthquakes only began in the early 1900s, so it is too early to categorically state that this is the case.[39]
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+ Most of the world's earthquakes (90%, and 81% of the largest) take place in the 40,000-kilometre-long (25,000 mi), horseshoe-shaped zone called the circum-Pacific seismic belt, known as the Pacific Ring of Fire, which for the most part bounds the Pacific Plate.[40][41] Massive earthquakes tend to occur along other plate boundaries too, such as along the Himalayan Mountains.[42]
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+ With the rapid growth of mega-cities such as Mexico City, Tokyo and Tehran in areas of high seismic risk, some seismologists are warning that a single quake may claim the lives of up to three million people.[43]
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+ While most earthquakes are caused by movement of the Earth's tectonic plates, human activity can also produce earthquakes. Activities both above ground and below may change the stresses and strains on the crust, including building reservoirs, extracting resources such as coal or oil, and injecting fluids underground for waste disposal or fracking.[44] Most of these earthquakes have small magnitudes. The 5.7 magnitude 2011 Oklahoma earthquake is thought to have been caused by disposing wastewater from oil production into injection wells,[45] and studies point to the state's oil industry as the cause of other earthquakes in the past century.[46] A Columbia University paper suggested that the 8.0 magnitude 2008 Sichuan earthquake was induced by loading from the Zipingpu Dam, though the link has not been conclusively proved.[47]
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+ The instrumental scales used to describe the size of an earthquake began with the Richter magnitude scale in the 1930s. It is a relatively simple measurement of an event's amplitude, and its use has become minimal in the 21st century. Seismic waves travel through the Earth's interior and can be recorded by seismometers at great distances. The surface wave magnitude was developed in the 1950s as a means to measure remote earthquakes and to improve the accuracy for larger events. The moment magnitude scale not only measures the amplitude of the shock but also takes into account the seismic moment (total rupture area, average slip of the fault, and rigidity of the rock). The Japan Meteorological Agency seismic intensity scale, the Medvedev–Sponheuer–Karnik scale, and the Mercalli intensity scale are based on the observed effects and are related to the intensity of shaking.
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+ Every tremor produces different types of seismic waves, which travel through rock with different velocities:
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+ Propagation velocity of the seismic waves through solid rock ranges from approx. 3 km/s (1.9 mi/s) up to 13 km/s (8.1 mi/s), depending on the density and elasticity of the medium. In the Earth's interior, the shock- or P-waves travel much faster than the S-waves (approx. relation 1.7:1). The differences in travel time from the epicenter to the observatory are a measure of the distance and can be used to image both sources of quakes and structures within the Earth. Also, the depth of the hypocenter can be computed roughly.
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+ In the upper crust, P-waves travel in the range 2–3 km (1.2–1.9 mi) per second (or lower) in soils and unconsolidated sediments, increasing to 3–6 km (1.9–3.7 mi) per second in solid rock. In the lower crust, they travel at about 6–7 km (3.7–4.3 mi) per second; the velocity increases within the deep mantle to about 13 km (8.1 mi) per second. The velocity of S-waves ranges from 2–3 km (1.2–1.9 mi) per second in light sediments and 4–5 km (2.5–3.1 mi) per second in the Earth's crust up to 7 km (4.3 mi) per second in the deep mantle. As a consequence, the first waves of a distant earthquake arrive at an observatory via the Earth's mantle.
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+ On average, the kilometer distance to the earthquake is the number of seconds between the P- and S-wave times 8.[48] Slight deviations are caused by inhomogeneities of subsurface structure. By such analyses of seismograms the Earth's core was located in 1913 by Beno Gutenberg.
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+ S-waves and later arriving surface waves do most of the damage compared to P-waves. P-waves squeeze and expand material in the same direction they are traveling, whereas S-waves shake the ground up and down and back and forth.[49]
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+ Earthquakes are not only categorized by their magnitude but also by the place where they occur. The world is divided into 754 Flinn–Engdahl regions (F-E regions), which are based on political and geographical boundaries as well as seismic activity. More active zones are divided into smaller F-E regions whereas less active zones belong to larger F-E regions.
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+ Standard reporting of earthquakes includes its magnitude, date and time of occurrence, geographic coordinates of its epicenter, depth of the epicenter, geographical region, distances to population centers, location uncertainty, a number of parameters that are included in USGS earthquake reports (number of stations reporting, number of observations, etc.), and a unique event ID.[50]
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+ Although relatively slow seismic waves have traditionally been used to detect earthquakes, scientists realized in 2016 that gravitational measurements could provide instantaneous detection of earthquakes, and confirmed this by analyzing gravitational records associated with the 2011 Tohoku-Oki ("Fukushima") earthquake.[51][52]
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+ The effects of earthquakes include, but are not limited to, the following:
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+ Shaking and ground rupture are the main effects created by earthquakes, principally resulting in more or less severe damage to buildings and other rigid structures. The severity of the local effects depends on the complex combination of the earthquake magnitude, the distance from the epicenter, and the local geological and geomorphological conditions, which may amplify or reduce wave propagation.[53] The ground-shaking is measured by ground acceleration.
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+ Specific local geological, geomorphological, and geostructural features can induce high levels of shaking on the ground surface even from low-intensity earthquakes. This effect is called site or local amplification. It is principally due to the transfer of the seismic motion from hard deep soils to soft superficial soils and to effects of seismic energy focalization owing to typical geometrical setting of the deposits.
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+ Ground rupture is a visible breaking and displacement of the Earth's surface along the trace of the fault, which may be of the order of several meters in the case of major earthquakes. Ground rupture is a major risk for large engineering structures such as dams, bridges, and nuclear power stations and requires careful mapping of existing faults to identify any that are likely to break the ground surface within the life of the structure.[54]
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+ Soil liquefaction occurs when, because of the shaking, water-saturated granular material (such as sand) temporarily loses its strength and transforms from a solid to a liquid. Soil liquefaction may cause rigid structures, like buildings and bridges, to tilt or sink into the liquefied deposits. For example, in the 1964 Alaska earthquake, soil liquefaction caused many buildings to sink into the ground, eventually collapsing upon themselves.[55]
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+ An earthquake may cause injury and loss of life, road and bridge damage, general property damage, and collapse or destabilization (potentially leading to future collapse) of buildings. The aftermath may bring disease, lack of basic necessities, mental consequences such as panic attacks, depression to survivors,[56] and higher insurance premiums.
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+ Earthquakes can produce slope instability leading to landslides, a major geological hazard. Landslide danger may persist while emergency personnel are attempting rescue.[57]
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+ Earthquakes can cause fires by damaging electrical power or gas lines. In the event of water mains rupturing and a loss of pressure, it may also become difficult to stop the spread of a fire once it has started. For example, more deaths in the 1906 San Francisco earthquake were caused by fire than by the earthquake itself.[58]
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+ Tsunamis are long-wavelength, long-period sea waves produced by the sudden or abrupt movement of large volumes of water—including when an earthquake occurs at sea. In the open ocean the distance between wave crests can surpass 100 kilometers (62 mi), and the wave periods can vary from five minutes to one hour. Such tsunamis travel 600–800 kilometers per hour (373–497 miles per hour), depending on water depth. Large waves produced by an earthquake or a submarine landslide can overrun nearby coastal areas in a matter of minutes. Tsunamis can also travel thousands of kilometers across open ocean and wreak destruction on far shores hours after the earthquake that generated them.[59]
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+ Ordinarily, subduction earthquakes under magnitude 7.5 do not cause tsunamis, although some instances of this have been recorded. Most destructive tsunamis are caused by earthquakes of magnitude 7.5 or more.[59]
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+ Floods may be secondary effects of earthquakes, if dams are damaged. Earthquakes may cause landslips to dam rivers, which collapse and cause floods.[60]
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+ The terrain below the Sarez Lake in Tajikistan is in danger of catastrophic flooding if the landslide dam formed by the earthquake, known as the Usoi Dam, were to fail during a future earthquake. Impact projections suggest the flood could affect roughly 5 million people.[61]
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+ One of the most devastating earthquakes in recorded history was the 1556 Shaanxi earthquake, which occurred on 23 January 1556 in Shaanxi province, China. More than 830,000 people died.[63] Most houses in the area were yaodongs—dwellings carved out of loess hillsides—and many victims were killed when these structures collapsed. The 1976 Tangshan earthquake, which killed between 240,000 and 655,000 people, was the deadliest of the 20th century.[64]
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+ The 1960 Chilean earthquake is the largest earthquake that has been measured on a seismograph, reaching 9.5 magnitude on 22 May 1960.[30][31] Its epicenter was near Cañete, Chile. The energy released was approximately twice that of the next most powerful earthquake, the Good Friday earthquake (March 27, 1964), which was centered in Prince William Sound, Alaska.[65][66] The ten largest recorded earthquakes have all been megathrust earthquakes; however, of these ten, only the 2004 Indian Ocean earthquake is simultaneously one of the deadliest earthquakes in history.
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+ Earthquakes that caused the greatest loss of life, while powerful, were deadly because of their proximity to either heavily populated areas or the ocean, where earthquakes often create tsunamis that can devastate communities thousands of kilometers away. Regions most at risk for great loss of life include those where earthquakes are relatively rare but powerful, and poor regions with lax, unenforced, or nonexistent seismic building codes.
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+ Earthquake prediction is a branch of the science of seismology concerned with the specification of the time, location, and magnitude of future earthquakes within stated limits.[67] Many methods have been developed for predicting the time and place in which earthquakes will occur. Despite considerable research efforts by seismologists, scientifically reproducible predictions cannot yet be made to a specific day or month.[68]
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+ While forecasting is usually considered to be a type of prediction, earthquake forecasting is often differentiated from earthquake prediction. Earthquake forecasting is concerned with the probabilistic assessment of general earthquake hazard, including the frequency and magnitude of damaging earthquakes in a given area over years or decades.[69] For well-understood faults the probability that a segment may rupture during the next few decades can be estimated.[70][71]
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+ Earthquake warning systems have been developed that can provide regional notification of an earthquake in progress, but before the ground surface has begun to move, potentially allowing people within the system's range to seek shelter before the earthquake's impact is felt.
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+ The objective of earthquake engineering is to foresee the impact of earthquakes on buildings and other structures and to design such structures to minimize the risk of damage. Existing structures can be modified by seismic retrofitting to improve their resistance to earthquakes. Earthquake insurance can provide building owners with financial protection against losses resulting from earthquakes Emergency management strategies can be employed by a government or organization to mitigate risks and prepare for consequences.
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+ Individuals can also take preparedness steps like securing water heaters and heavy items that could injure someone, locating shutoffs for utilities, and being educated about what to do when shaking starts. For areas near large bodies of water, earthquake preparedness encompasses the possibility of a tsunami caused by a large quake.
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+ From the lifetime of the Greek philosopher Anaxagoras in the 5th century BCE to the 14th century CE, earthquakes were usually attributed to "air (vapors) in the cavities of the Earth."[72] Thales of Miletus (625–547 BCE) was the only documented person who believed that earthquakes were caused by tension between the earth and water.[72] Other theories existed, including the Greek philosopher Anaxamines' (585–526 BCE) beliefs that short incline episodes of dryness and wetness caused seismic activity. The Greek philosopher Democritus (460–371 BCE) blamed water in general for earthquakes.[72] Pliny the Elder called earthquakes "underground thunderstorms".[72]
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+ In recent studies, geologists claim that global warming is one of the reasons for increased seismic activity. According to these studies, melting glaciers and rising sea levels disturb the balance of pressure on Earth's tectonic plates, thus causing an increase in the frequency and intensity of earthquakes.[73]
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+ In Norse mythology, earthquakes were explained as the violent struggling of the god Loki. When Loki, god of mischief and strife, murdered Baldr, god of beauty and light, he was punished by being bound in a cave with a poisonous serpent placed above his head dripping venom. Loki's wife Sigyn stood by him with a bowl to catch the poison, but whenever she had to empty the bowl the poison dripped on Loki's face, forcing him to jerk his head away and thrash against his bonds, which caused the earth to tremble.[74]
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+ In Greek mythology, Poseidon was the cause and god of earthquakes. When he was in a bad mood, he struck the ground with a trident, causing earthquakes and other calamities. He also used earthquakes to punish and inflict fear upon people as revenge.[75]
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+ In Japanese mythology, Namazu (鯰) is a giant catfish who causes earthquakes. Namazu lives in the mud beneath the earth, and is guarded by the god Kashima who restrains the fish with a stone. When Kashima lets his guard fall, Namazu thrashes about, causing violent earthquakes.[76]
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+ In modern popular culture, the portrayal of earthquakes is shaped by the memory of great cities laid waste, such as Kobe in 1995 or San Francisco in 1906.[77] Fictional earthquakes tend to strike suddenly and without warning.[77] For this reason, stories about earthquakes generally begin with the disaster and focus on its immediate aftermath, as in Short Walk to Daylight (1972), The Ragged Edge (1968) or Aftershock: Earthquake in New York (1999).[77] A notable example is Heinrich von Kleist's classic novella, The Earthquake in Chile, which describes the destruction of Santiago in 1647. Haruki Murakami's short fiction collection After the Quake depicts the consequences of the Kobe earthquake of 1995.
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+ The most popular single earthquake in fiction is the hypothetical "Big One" expected of California's San Andreas Fault someday, as depicted in the novels Richter 10 (1996), Goodbye California (1977), 2012 (2009) and San Andreas (2015) among other works.[77] Jacob M. Appel's widely anthologized short story, A Comparative Seismology, features a con artist who convinces an elderly woman that an apocalyptic earthquake is imminent.[78]
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+ Contemporary depictions of earthquakes in film are variable in the manner in which they reflect human psychological reactions to the actual trauma that can be caused to directly afflicted families and their loved ones.[79] Disaster mental health response research emphasizes the need to be aware of the different roles of loss of family and key community members, loss of home and familiar surroundings, loss of essential supplies and services to maintain survival.[80][81] Particularly for children, the clear availability of caregiving adults who are able to protect, nourish, and clothe them in the aftermath of the earthquake, and to help them make sense of what has befallen them has been shown even more important to their emotional and physical health than the simple giving of provisions.[82] As was observed after other disasters involving destruction and loss of life and their media depictions, recently observed in the 2010 Haiti earthquake, it is also important not to pathologize the reactions to loss and displacement or disruption of governmental administration and services, but rather to validate these reactions, to support constructive problem-solving and reflection as to how one might improve the conditions of those affected.[83]
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+ Tyrannosaurus[nb 1] is a genus of coelurosaurian theropod dinosaurs. The species Tyrannosaurus rex (rex meaning "king" in Latin), often called T. rex or colloquially T-Rex, is one of the most well-represented of the large theropods. Tyrannosaurus lived throughout what is now western North America, on what was then an island continent known as Laramidia. Tyrannosaurus had a much wider range than other tyrannosaurids. Fossils are found in a variety of rock formations dating to the Maastrichtian age of the upper Cretaceous period, 68 to 66 million years ago. It was the last known member of the tyrannosaurids, and among the last non-avian dinosaurs to exist before the Cretaceous–Paleogene extinction event.
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+ Like other tyrannosaurids, Tyrannosaurus was a bipedal carnivore with a massive skull balanced by a long, heavy tail. Relative to its large and powerful hind limbs, Tyrannosaurus forelimbs were short but unusually powerful for their size and had two clawed digits. The most complete specimen measures up to 12.3 meters (40 feet) in length though T. rex could grow to lengths of over 12.3 m (40 ft), up to 3.66 m (12 ft) tall at the hips, and according to most modern estimates 8.4 metric tons (9.3 short tons) to 14 metric tons (15.4 short tons) in weight. Although other theropods rivaled or exceeded Tyrannosaurus rex in size, it is still among the largest known land predators and is estimated to have exerted the strongest bite force among all terrestrial animals. By far the largest carnivore in its environment, Tyrannosaurus rex was most likely an apex predator, preying upon hadrosaurs, armored herbivores like ceratopsians and ankylosaurs, and possibly sauropods. Some experts have suggested the dinosaur was primarily a scavenger. The question of whether Tyrannosaurus was an apex predator or a pure scavenger was among the longest debates in paleontology. Most paleontologists today accept that Tyrannosaurus was both an active predator and a scavenger.
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+ Specimens of Tyrannosaurus rex include some that are nearly complete skeletons. Soft tissue and proteins have been reported in at least one of these specimens. The abundance of fossil material has allowed significant research into many aspects of its biology, including its life history and biomechanics. The feeding habits, physiology and potential speed of Tyrannosaurus rex are a few subjects of debate. Its taxonomy is also controversial, as some scientists consider Tarbosaurus bataar from Asia to be a second Tyrannosaurus species while others maintain Tarbosaurus is a separate genus. Several other genera of North American tyrannosaurids have also been synonymized with Tyrannosaurus.
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+ As the archetypal theropod, Tyrannosaurus has been one of the best-known dinosaurs since the early 20th century, and has been featured in film, advertising, postal stamps, and many other media.
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+ Teeth from what is now documented as a Tyrannosaurus rex were found in 1874 by Arthur Lakes near Golden, Colorado. In the early 1890s, John Bell Hatcher collected postcranial elements in eastern Wyoming. The fossils were believed to be from the large species Ornithomimus grandis (now Deinodon) but are now considered T. rex remains.[2]
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+ In 1892, Edward Drinker Cope found two vertebral fragments of large dinosaur. Cope believed the fragments belonged to an "agathaumid" (ceratopsid) dinosaur, and named them Manospondylus gigas, meaning "giant porous vertebra", in reference to the numerous openings for blood vessels he found in the bone.[2] The M. gigas remains were, in 1907, identified by Hatcher as those of a theropod rather than a ceratopsid.[3]
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+ Henry Fairfield Osborn recognized the similarity between Manospondylus gigas and T. rex as early as 1917, by which time the second vertebra had been lost. Owing to the fragmentary nature of the Manospondylus vertebrae, Osborn did not synonymize the two genera, instead considering the older genus indeterminate.[4] In June 2000, the Black Hills Institute found around 10% of a Tyrannosaurus skeleton (BHI 6248) at a site that might have been the original M. gigas locality.[5]
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+ Barnum Brown, assistant curator of the American Museum of Natural History, found the first partial skeleton of T. rex in eastern Wyoming in 1900. Brown found another partial skeleton in the Hell Creek Formation in Montana in 1902, comprising approximately 34 fossilized bones.[6] Writing at the time Brown said "Quarry No. 1 contains the femur, pubes, humerus, three vertebrae and two undetermined bones of a large Carnivorous Dinosaur not described by Marsh.... I have never seen anything like it from the Cretaceous".[7] Henry Fairfield Osborn, president of the American Museum of Natural History, named the second skeleton T. rex in 1905. The generic name is derived from the Greek words τύραννος (tyrannos, meaning "tyrant") and σαῦρος (sauros, meaning "lizard"). Osborn used the Latin word rex, meaning "king", for the specific name. The full binomial therefore translates to "tyrant lizard the king" or "King Tyrant Lizard", emphasizing the animal's size and perceived dominance over other species of the time.[6]
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+ Osborn named the other specimen Dynamosaurus imperiosus in a paper in 1905.[6] In 1906, Osborn recognized that the two skeletons were from the same species and selected Tyrannosaurus as the preferred name.[8] The original Dynamosaurus material resides in the collections of the Natural History Museum, London.[9] In 1941, the T. rex type specimen was sold to the Carnegie Museum of Natural History in Pittsburgh, Pennsylvania, for $7,000.[7] Dynamosaurus would later be honored by the 2018 description of another species of tyrannosaurid by Andrew McDonald and colleagues, Dynamoterror dynastes, whose name was chosen in reference to the 1905 name, as it had been a "childhood favorite" of McDonald's.[10]
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+ From the 1910s through the end of the 1950s, Barnum's discoveries remained the only specimens of Tyrannosaurus, as the Great Depression and wars kept many paleontologists out of the field.[5]
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+ Beginning in the 1960s, there was renewed interest in Tyrannosaurus, resulting in recovery of 42 skeletons (5-80% complete by bone count) from Western North America.[5] In 1967, Dr. William MacMannis located and recovered the skeleton named "MOR 008", which is 15% complete by bone count and has a reconstructed skull displayed at the Museum of the Rockies. The 1990s saw numerous discoveries, with nearly twice as many finds as in all previous years, including two of the most complete skeletons found to date: Sue and Stan.[5]
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+ Sue Hendrickson, an amateur paleontologist, discovered the most complete (approximately 85%) and largest Tyrannosaurus skeleton in the Hell Creek Formation on August 12, 1990. The specimen Sue, named after the discoverer, was the object of a legal battle over its ownership. In 1997, the litigation was settled in favor of Maurice Williams, the original land owner. The fossil collection was purchased by the Field Museum of Natural History at auction for $7.6 million, making it the most expensive dinosaur skeleton to date. From 1998 to 1999, Field Museum of Natural History staff spent over 25,000 hours taking the rock off the bones.[11] The bones were then shipped to New Jersey where the mount was constructed, then shipped back to Chicago for the final assembly. The mounted skeleton opened to the public on May 17, 2000 in the Field Museum of Natural History. A study of this specimen's fossilized bones showed that Sue reached full size at age 19 and died at the age of 28, the longest estimated life of any tyrannosaur known.[12]
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+ Another Tyrannosaurus, nicknamed Stan (BHI 3033), in honor of amateur paleontologist Stan Sacrison, was recovered from the Hell Creek Formation in 1992. Stan is the second most complete skeleton found, with 199 bones recovered representing 70% of the total.[13] This tyrannosaur also had many bone pathologies, including broken and healed ribs, a broken (and healed) neck, and a substantial hole in the back of its head, about the size of a Tyrannosaurus tooth.[14]
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+ In 1998, Bucky Derflinger noticed one of Bucky's toes exposed above ground, making Derflinger, who was 20 years old at the time, the youngest person to discover a Tyrannosaurus. The specimen was a young adult, 3.0 metres (10 ft) tall and 11 metres (35 ft) long. Bucky is the first Tyrannosaurus to be found that preserved a furcula (wishbone). Bucky is permanently displayed at The Children's Museum of Indianapolis.[15]
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+ In the summer of 2000, crews organized by Jack Horner discovered five Tyrannosaurus skeletons near the Fort Peck Reservoir.[16] In 2001, a 50% complete skeleton of a juvenile Tyrannosaurus was discovered in the Hell Creek Formation by a crew from the Burpee Museum of Natural History. Dubbed Jane (BMRP 2002.4.1), the find was thought to be the first known skeleton of a pygmy tyrannosaurid, Nanotyrannus, but subsequent research revealed that it is more likely a juvenile Tyrannosaurus, and the most complete juvenile example known;[17] Jane is exhibited at the Burpee Museum of Natural History.[18] In 2002, a skeleton named Wyrex, discovered by amateur collectors Dan Wells and Don Wyrick, had 114 bones and was 38% complete. The dig was concluded over 3 weeks in 2004 by the Black Hills Institute with the first live online Tyrannosaurus excavation providing daily reports, photos, and video.[5]
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+ In 2006, Montana State University revealed that it possessed the largest Tyrannosaurus skull yet discovered (from a specimen named MOR 008), measuring 5 feet (152 cm) long.[19] Subsequent comparisons indicated that the longest head was 136.5 centimetres (53.7 in) (from specimen LACM 23844) and the widest head was 90.2 centimetres (35.5 in) (from Sue).[20]
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+ Two isolated fossilized footprints have been tentatively assigned to T. rex. The first was discovered at Philmont Scout Ranch, New Mexico, in 1983 by American geologist Charles Pillmore. Originally thought to belong to a hadrosaurid, examination of the footprint revealed a large 'heel' unknown in ornithopod dinosaur tracks, and traces of what may have been a hallux, the dewclaw-like fourth digit of the tyrannosaur foot. The footprint was published as the ichnogenus Tyrannosauripus pillmorei in 1994, by Martin Lockley and Adrian Hunt. Lockley and Hunt suggested that it was very likely the track was made by a T. rex, which would make it the first known footprint from this species. The track was made in what was once a vegetated wetland mud flat. It measures 83 centimeters (33 in) long by 71 centimeters (28 in) wide.[21]
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+ A second footprint that may have been made by a Tyrannosaurus was first reported in 2007 by British paleontologist Phil Manning, from the Hell Creek Formation of Montana. This second track measures 72 centimeters (28 in) long, shorter than the track described by Lockley and Hunt. Whether or not the track was made by Tyrannosaurus is unclear, though Tyrannosaurus and Nanotyrannus are the only large theropods known to have existed in the Hell Creek Formation.[22][23]
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+ A set of footprints in Glenrock, Wyoming dating to the Maastrichtian stage of the Late Cretaceous and hailing from the Lance Formation were described by Scott Persons, Phil Currie and colleagues in 2016, and are believed to belong to either a juvenile T. rex or the dubious tyrannosaurid Nanotyrannus lancensis. From measurements and based on the positions of the footprints, the animal was believed to be traveling at a walking speed of around 2.8 to 5 miles per hour and was estimated to have a hip height of 1.56 m (5.1 ft) to 2.06 m (6.8 ft).[24][25][26] A follow-up paper appeared in 2017, increasing the speed estimations by 50-80%.[27]
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+ T. rex was one of the largest land carnivores of all time. One of the largest and the most complete specimens, nicknamed Sue (FMNH PR2081), is located at the Field Museum of Natural History. Sue measured 12.3–12.8 meters (40–42 ft) long,[28][29] was 3.66 meters (12 ft) tall at the hips,[30] and according to the most recent studies, using a variety of techniques, estimated to have weighed between 8.4 metric tons (9.3 short tons) to 14 metric tons (15.4 short tons).[29][31] A specimen nicknamed Scotty (RSM P2523.8), located at the Royal Saskatchewan Museum, is reported to measure 13 m (43 ft) in length. Using a mass estimation technique that extrapolates from the circumference of the femur, Scotty was estimated as the largest known specimen at 8.8 metric tons (9.7 short tons) in weight.[32][33]
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+ Not every adult Tyrannosaurus specimen recovered is as big. Historically average adult mass estimates have varied widely over the years, from as low as 4.5 metric tons (5.0 short tons),[34][35] to more than 7.2 metric tons (7.9 short tons),[36] with most modern estimates ranging between 5.4 metric tons (6.0 short tons) and 8.0 metric tons (8.8 short tons).[29][37][38][39][40]
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+ The largest known T. rex skull is 1.52 meters (5 ft) in length.[30] Large fenestrae (openings) in the skull reduced weight, as in all carnivorous theropods. In other respects Tyrannosaurus's skull was significantly different from those of large non-tyrannosaurid theropods. It was extremely wide at the rear but had a narrow snout, allowing unusually good binocular vision.[41][42] The skull bones were massive and the nasals and some other bones were fused, preventing movement between them; but many were pneumatized (contained a "honeycomb" of tiny air spaces) and thus lighter. These and other skull-strengthening features are part of the tyrannosaurid trend towards an increasingly powerful bite, which easily surpassed that of all non-tyrannosaurids.[43][44][45] The tip of the upper jaw was U-shaped (most non-tyrannosauroid carnivores had V-shaped upper jaws), which increased the amount of tissue and bone a tyrannosaur could rip out with one bite, although it also increased the stresses on the front teeth.[46]
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+ The teeth of T. rex displayed marked heterodonty (differences in shape).[47][48] The premaxillary teeth, four per side at the front of the upper jaw, were closely packed, D-shaped in cross-section, had reinforcing ridges on the rear surface, were incisiform (their tips were chisel-like blades) and curved backwards. The D-shaped cross-section, reinforcing ridges and backwards curve reduced the risk that the teeth would snap when Tyrannosaurus bit and pulled. The remaining teeth were robust, like "lethal bananas" rather than daggers, more widely spaced and also had reinforcing ridges.[49] Those in the upper jaw, twelve per side in mature individuals,[47] were larger than their counterparts of the lower jaw, except at the rear. The largest found so far is estimated to have been 30.5 centimeters (12 in) long including the root when the animal was alive, making it the largest tooth of any carnivorous dinosaur yet found.[50] The lower jaw was robust. Its front dentary bone bore thirteen teeth. Behind the tooth row, the lower jaw became notably taller.[47] The upper and lower jaws of Tyrannosaurus, like those of many dinosaurs, possessed numerous foramina, or small holes in the bone. Various functions have been proposed for these foramina, such as a crocodile-like sensory system[51] or evidence of extra-oral structures such as scales or potentially lips.[52][53]
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+ The vertebral column of Tyrannosaurus consisted of ten neck vertebrae, thirteen back vertebrae and five sacral vertebrae. The number of tail vertebrae is unknown and could well have varied between individuals but probably numbered at least forty. Sue was mounted with forty-seven of such caudal vertebrae.[47] The neck of T. rex formed a natural S-shaped curve like that of other theropods. Compared to these, it was exceptionally short, deep and muscular to support the massive head. The second vertebra, the axis, was especially short. The remaining neck vertebrae were weakly opisthocoelous, i.e. with a convex front of the vertebral body and a concave rear. The vertebral bodies had single pleurocoels, pneumatic depressions created by air sacs, on their sides.[47] The vertebral bodies of the torso were robust but with a narrow waist. Their undersides were keeled. The front sides were concave with a deep vertical trough. They had large pleurocoels. Their neural spines had very rough front and rear sides for the attachment of strong tendons. The sacral vertebrae were fused to each other, both in their vertebral bodies and neural spines. They were pneumatized. They were connected to the pelvis by transverse processes and sacral ribs. The tail was heavy and moderately long, in order to balance the massive head and torso and to provide space for massive locomotor muscles that attached to the thighbones. The thirteenth tail vertebra formed the transition point between the deep tail base and the middle tail that was stiffened by rather long front articulation processes. The underside of the trunk was covered by eighteen or nineteen pairs of segmented belly ribs.[47]
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+ The shoulder girdle was longer than the entire forelimb. The shoulder blade had a narrow shaft but was exceptionally expanded at its upper end. It connected via a long forward protrusion to the coracoid, which was rounded. Both shoulder blades were connected by a small furcula. The paired breast bones possibly were made of cartilage only.[47]
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+ The forelimb or arm was very short. The upper arm bone, the humerus, was short but robust. It had a narrow upper end with an exceptionally rounded head. The lower arm bones, the ulna and radius, were straight elements, much shorter than the humerus. The second metacarpal was longer and wider than the first, whereas normally in theropods the opposite is true. The forelimbs had only two clawed fingers,[47] along with an additional splint-like small third metacarpal representing the remnant of a third digit.[54]
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+ The pelvis was a large structure. Its upper bone, the ilium, was both very long and high, providing an extensive attachment area for hindlimb muscles. The front pubic bone ended in an enormous pubic boot, longer than the entire shaft of the element. The rear ischium was slender and straight, pointing obliquely to behind and below.[47]
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+ In contrast to the arms, the hindlimbs were among the longest in proportion to body size of any theropod. In the foot, the metatarsus was "arctometatarsalian", meaning that the part of the third metatarsal near the ankle was pinched. The third metatarsal was also exceptionally sinuous.[47] Compensating for the immense bulk of the animal, many bones throughout the skeleton were hollowed, reducing its weight without significant loss of strength.[47]
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+ The discovery of feathered dinosaurs led to debate regarding whether, and to what extent, Tyrannosaurus might have been feathered.[55][56] Filamentous structures, which are commonly recognized as the precursors of feathers, have been reported in the small-bodied, basal tyrannosauroid Dilong paradoxus from the Early Cretaceous Yixian Formation of China in 2004.[57] Because integumentary impressions of larger tyrannosauroids known at that time showed evidence of scales, the researchers who studied Dilong speculated that insulating feathers might have been lost by larger species due to their smaller surface-to-volume ratio.[57] The subsequent discovery of the giant species Yutyrannus huali, also from the Yixian, showed that even some large tyrannosauroids had feathers covering much of their bodies, casting doubt on the hypothesis that they were a size-related feature.[58] A 2017 study reviewed known skin impressions of tyrannosaurids, including those of a Tyrannosaurus specimen nicknamed "Wyrex" (BHI 6230) which preserves patches of mosaic scales on the tail, hip, and neck.[5] The study concluded that feather covering of large tyrannosaurids such as Tyrannosaurus was, if present, limited to the upper side of the trunk.[55]
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+ A conference abstract published in 2016 posited that theropods such as Tyrannosaurus had their upper teeth covered in lips, instead of bare teeth as seen in crocodilians. This was based on the presence of enamel, which according to the study needs to remain hydrated, an issue not faced by aquatic animals like crocodilians.[53] A 2017 analytical study proposed that tyrannosaurids had large, flat scales on their snouts instead of lips.[51][59]
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+ Tyrannosaurus is the type genus of the superfamily Tyrannosauroidea, the family Tyrannosauridae, and the subfamily Tyrannosaurinae; in other words it is the standard by which paleontologists decide whether to include other species in the same group. Other members of the tyrannosaurine subfamily include the North American Daspletosaurus and the Asian Tarbosaurus,[17][60] both of which have occasionally been synonymized with Tyrannosaurus.[61] Tyrannosaurids were once commonly thought to be descendants of earlier large theropods such as megalosaurs and carnosaurs, although more recently they were reclassified with the generally smaller coelurosaurs.[46]
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+ In 1955, Soviet paleontologist Evgeny Maleev named a new species, Tyrannosaurus bataar, from Mongolia.[62] By 1965, this species had been renamed Tarbosaurus bataar.[63] Despite the renaming, many phylogenetic analyses have found Tarbosaurus bataar to be the sister taxon of T. rex,[60] and it has often been considered an Asian species of Tyrannosaurus.[46][64][65] The discovery of the tyrannosaurid Lythronax further indicates that Tarbosaurus and Tyrannosaurus are closely related, forming a clade with fellow Asian tyrannosaurid Zhuchengtyrannus, with Lythronax being their sister taxon.[66][67] A further study from 2016 by Steve Brusatte, Thomas Carr and colleagues, also indicates that Tyrannosaurus may have been an immigrant from Asia, as well as a possible descendant of Tarbosaurus.[68]
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+ In 2001, various tyrannosaurid teeth and a metatarsal unearthed in a quarry near Zhucheng, China were assigned by Chinese paleontologist Hu Chengzhi to the newly erected Tyrannosaurus zhuchengensis. However, in a nearby site, a right maxilla and left jawbone were assigned to the newly erected tyrannosaurid genus Zhuchengtyrannus in 2011, and it is possible T. zhuchengensis is synonymous with Zhuchengtyrannus. In any case, T. zhuchengensis is considered to be a nomen dubium as the holotype lacks diagnostic features below the level Tyrannosaurinae.[69]
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+ Below is the cladogram of Tyrannosauridae based on the phylogenetic analysis conducted by Loewen and colleagues in 2013.[66]
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+ Gorgosaurus libratus
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+
75
+ Albertosaurus sarcophagus
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+
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+ Dinosaur Park tyrannosaurid
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+
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+ Daspletosaurus torosus
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+
81
+ Two Medicine tyrannosaurid
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+
83
+ Teratophoneus curriei
84
+
85
+ Bistahieversor sealeyi
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+
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+ Lythronax argestes
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+
89
+ Tyrannosaurus rex
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+
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+ Tarbosaurus bataar
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+
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+ Zhuchengtyrannus magnus
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+
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+ Other tyrannosaurid fossils found in the same formations as T. rex were originally classified as separate taxa, including Aublysodon and Albertosaurus megagracilis,[61] the latter being named Dinotyrannus megagracilis in 1995.[70] These fossils are now universally considered to belong to juvenile T. rex.[71] A small but nearly complete skull from Montana, 60 centimeters (2.0 ft) long, might be an exception. This skull, CMNH 7541, was originally classified as a species of Gorgosaurus (G. lancensis) by Charles W. Gilmore in 1946.[72] In 1988, the specimen was re-described by Robert T. Bakker, Phil Currie, and Michael Williams, then the curator of paleontology at the Cleveland Museum of Natural History, where the original specimen was housed and is now on display. Their initial research indicated that the skull bones were fused, and that it therefore represented an adult specimen. In light of this, Bakker and colleagues assigned the skull to a new genus named Nanotyrannus (meaning "dwarf tyrant", for its apparently small adult size). The specimen is estimated to have been around 5.2 metres (17 ft) long when it died.[73] However, In 1999, a detailed analysis by Thomas Carr revealed the specimen to be a juvenile, leading Carr and many other paleontologists to consider it a juvenile T. rex individual.[74][75]
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+ In 2001, a more complete juvenile tyrannosaur (nicknamed "Jane", catalog number BMRP 2002.4.1), belonging to the same species as the original Nanotyrannus specimen, was uncovered. This discovery prompted a conference on tyrannosaurs focused on the issues of Nanotyrannus validity at the Burpee Museum of Natural History in 2005. Several paleontologists who had previously published opinions that N. lancensis was a valid species, including Currie and Williams, saw the discovery of "Jane" as a confirmation that Nanotyrannus was, in fact, a juvenile T. rex.[76][77][78] Peter Larson continued to support the hypothesis that N. lancensis was a separate but closely related species, based on skull features such as two more teeth in both jaws than T. rex; as well as proportionately larger hands with phalanges on the third metacarpal and different wishbone anatomy in an undescribed specimen. He also argued that Stygivenator, generally considered to be a juvenile T. rex, may be a younger Nanotyrannus specimen.[79][80] Later research revealed that other tyrannosaurids such as Gorgosaurus also experienced reduction in tooth count during growth,[74] and given the disparity in tooth count between individuals of the same age group in this genus and Tyrannosaurus, this feature may also be due to individual variation.[75] In 2013, Carr noted that all of the differences claimed to support Nanotyrannus have turned out to be individually or ontogenetically variable features or products of distortion of the bones.[81]
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+ In 2016, analysis of limb proportions by Persons and Currie suggested Nanotyrannus specimens to have differing cursoriality levels, potentially separating it from T. rex.[82] However, paleontologist Manabu Sakomoto has commented that this conclusion may be impacted by low sample size, and the discrepancy does not necessarily reflect taxonomic distinction.[83] In 2016, Joshua Schmerge argued for Nanotyrannus' validity based on skull features, including a dentary groove in BMRP 2002.4.1's skull. According to Schmerge, as that feature is absent in T. rex and found only in Dryptosaurus and albertosaurines, this suggests Nanotyrannus is a distinct taxon within the Albertosaurinae.[84] The same year, Carr and colleagues noted that this was not sufficient enough to clarify Nanotyrannus' validity or classification, being a common and ontogenetically variable feature among tyrannosauroids.[85]
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+ A 2020 study by Holly Woodward and colleagues showed the specimens referred to Nanotyrannus were all ontogenetically immature and found it probable that these specimens belonged to T. rex.[86] The same year, Carr published a paper on T. rex's growth history, finding that CMNH 7541 fit within the expected ontogenetic variation of the taxon and displayed juvenile characteristics found in other specimens. It was classified as a juvenile, under 13 years old with a skull less than 80 cm (31 in). No significant sexual or phylogenetic variation was discernible among any of the 44 specimens studied, with Carr stating that characters of potential phylogenetic importance decrease throughout age at the same rate as growth occurs.[87] Discussing the paper's results, Carr described how all "Nanotyrannus" specimens formed a continual growth transition between the smallest juveniles and the subadults, unlike what would be expected if it were a distinct taxon where the specimens would group to the exclusion of Tyrannosaurus. Carr concluded that "the 'nanomorphs' are not all that similar to each other and instead form an important bridge in the growth series of T. rex that captures the beginnings of the profound change from the shallow skull of juveniles to the deep skull that is seen in fully-developed adults."[88]
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+ The identification of several specimens as juvenile T. rex has allowed scientists to document ontogenetic changes in the species, estimate the lifespan, and determine how quickly the animals would have grown. The smallest known individual (LACM 28471, the "Jordan theropod") is estimated to have weighed only 30 kg (66 lb), while the largest, such as FMNH PR2081 (Sue) most likely weighed about 5,650 kg (12,460 lb). Histologic analysis of T. rex bones showed LACM 28471 had aged only 2 years when it died, while Sue was 28 years old, an age which may have been close to the maximum for the species.[37]
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+ Histology has also allowed the age of other specimens to be determined. Growth curves can be developed when the ages of different specimens are plotted on a graph along with their mass. A T. rex growth curve is S-shaped, with juveniles remaining under 1,800 kg (4,000 lb) until approximately 14 years of age, when body size began to increase dramatically. During this rapid growth phase, a young T. rex would gain an average of 600 kg (1,300 lb) a year for the next four years. At 18 years of age, the curve plateaus again, indicating that growth slowed dramatically. For example, only 600 kg (1,300 lb) separated the 28-year-old Sue from a 22-year-old Canadian specimen (RTMP 81.12.1).[37] A 2004 histological study performed by different workers corroborates these results, finding that rapid growth began to slow at around 16 years of age.[89]
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+ A study by Hutchinson and colleagues in 2011 corroborated the previous estimation methods in general, but their estimation of peak growth rates is significantly higher; it found that the "maximum growth rates for T. rex during the exponential stage are 1790 kg/year".[29] Although these results were much higher than previous estimations, the authors noted that these results significantly lowered the great difference between its actual growth rate and the one which would be expected of an animal of its size.[29] The sudden change in growth rate at the end of the growth spurt may indicate physical maturity, a hypothesis which is supported by the discovery of medullary tissue in the femur of a 16 to 20-year-old T. rex from Montana (MOR 1125, also known as B-rex). Medullary tissue is found only in female birds during ovulation, indicating that B-rex was of reproductive age.[90] Further study indicates an age of 18 for this specimen.[91] In 2016, it was finally confirmed by Mary Higby Schweitzer and Lindsay Zanno and colleagues that the soft tissue within the femur of MOR 1125 was medullary tissue. This also confirmed the identity of the specimen as a female. The discovery of medullary bone tissue within Tyrannosaurus may prove valuable in determining the sex of other dinosaur species in future examinations, as the chemical makeup of medullary tissue is unmistakable.[92] Other tyrannosaurids exhibit extremely similar growth curves, although with lower growth rates corresponding to their lower adult sizes.[93]
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+ An additional study published in 2020 by Woodward and colleagues, for the journal Science Advances indicates that during their growth from juvenile to adult, Tyrannosaurus was capable of slowing down its growth to counter environmental factors such as lack of food. The study, focusing on two juvenile specimens between 13 and 15 years old housed at the Burpee Museum in Illinois, indicates that the rate of maturation for Tyrannosaurus was dependent on resource abundance. This study also indicates that in such changing environments, Tyrannosaurus was particularly well-suited to an environment that shifted yearly in regards to resource abundance, hinting that other midsize predators might have had difficulty surviving in such harsh conditions and explaining the niche partitioning between juvenile and adult tyrannosaurs. The study further indicates that Tyrannosaurus and the dubious genus Nanotyrannus are synonymous, due to analysis of the growth rings in the bones of the two specimens studied.[94][95]
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+ Over half of the known T. rex specimens appear to have died within six years of reaching sexual maturity, a pattern which is also seen in other tyrannosaurs and in some large, long-lived birds and mammals today. These species are characterized by high infant mortality rates, followed by relatively low mortality among juveniles. Mortality increases again following sexual maturity, partly due to the stresses of reproduction. One study suggests that the rarity of juvenile T. rex fossils is due in part to low juvenile mortality rates; the animals were not dying in large numbers at these ages, and so were not often fossilized. This rarity may also be due to the incompleteness of the fossil record or to the bias of fossil collectors towards larger, more spectacular specimens.[93] In a 2013 lecture, Thomas Holtz Jr. suggested that dinosaurs "lived fast and died young" because they reproduced quickly whereas mammals have long life spans because they take longer to reproduce.[96] Gregory S. Paul also writes that Tyrannosaurus reproduced quickly and died young, but attributes their short life spans to the dangerous lives they lived.[97]
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+ As the number of known specimens increased, scientists began to analyze the variation between individuals and discovered what appeared to be two distinct body types, or morphs, similar to some other theropod species. As one of these morphs was more solidly built, it was termed the 'robust' morph while the other was termed 'gracile'. Several morphological differences associated with the two morphs were used to analyze sexual dimorphism in T. rex, with the 'robust' morph usually suggested to be female. For example, the pelvis of several 'robust' specimens seemed to be wider, perhaps to allow the passage of eggs.[98] It was also thought that the 'robust' morphology correlated with a reduced chevron on the first tail vertebra, also ostensibly to allow eggs to pass out of the reproductive tract, as had been erroneously reported for crocodiles.[99]
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+ In recent years, evidence for sexual dimorphism has been weakened. A 2005 study reported that previous claims of sexual dimorphism in crocodile chevron anatomy were in error, casting doubt on the existence of similar dimorphism between T. rex sexes.[100] A full-sized chevron was discovered on the first tail vertebra of Sue, an extremely robust individual, indicating that this feature could not be used to differentiate the two morphs anyway. As T. rex specimens have been found from Saskatchewan to New Mexico, differences between individuals may be indicative of geographic variation rather than sexual dimorphism. The differences could also be age-related, with 'robust' individuals being older animals.[47]
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+ Only a single T. rex specimen has been conclusively shown to belong to a specific sex. Examination of B-rex demonstrated the preservation of soft tissue within several bones. Some of this tissue has been identified as a medullary tissue, a specialized tissue grown only in modern birds as a source of calcium for the production of eggshell during ovulation. As only female birds lay eggs, medullary tissue is only found naturally in females, although males are capable of producing it when injected with female reproductive hormones like estrogen. This strongly suggests that B-rex was female, and that she died during ovulation.[90] Recent research has shown that medullary tissue is never found in crocodiles, which are thought to be the closest living relatives of dinosaurs, aside from birds. The shared presence of medullary tissue in birds and theropod dinosaurs is further evidence of the close evolutionary relationship between the two.[101]
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+ Like many bipedal dinosaurs, T. rex was historically depicted as a 'living tripod', with the body at 45 degrees or less from the vertical and the tail dragging along the ground, similar to a kangaroo. This concept dates from Joseph Leidy's 1865 reconstruction of Hadrosaurus, the first to depict a dinosaur in a bipedal posture.[102] In 1915, convinced that the creature stood upright, Henry Fairfield Osborn, former president of the American Museum of Natural History, further reinforced the notion in unveiling the first complete T. rex skeleton arranged this way. It stood in an upright pose for 77 years, until it was dismantled in 1992.[103]
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+ By 1970, scientists realized this pose was incorrect and could not have been maintained by a living animal, as it would have resulted in the dislocation or weakening of several joints, including the hips and the articulation between the head and the spinal column.[104] The inaccurate AMNH mount inspired similar depictions in many films and paintings (such as Rudolph Zallinger's famous mural The Age of Reptiles in Yale University's Peabody Museum of Natural History)[105] until the 1990s, when films such as Jurassic Park introduced a more accurate posture to the general public.[106] Modern representations in museums, art, and film show T. rex with its body approximately parallel to the ground with the tail extended behind the body to balance the head.[107]
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+ To sit down, Tyrannosaurus may have settled its weight backwards and rested its weight on a pubic boot, the wide expansion at the end of the pubis in some dinosaurs. With its weight rested on the pelvis, it may have been free to move the hindlimbs. Getting back up again might have involved some stabilization from the diminutive forelimbs.[108][104] The latter known as Newman's pushup theory has been debated. Nonetheless, Tyrannosaurus was probably able to get up if it fell, which only would have required placing the limbs below the center of gravity, with the tail as an effective counterbalance.[109]
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+ When T. rex was first discovered, the humerus was the only element of the forelimb known.[6] For the initial mounted skeleton as seen by the public in 1915, Osborn substituted longer, three-fingered forelimbs like those of Allosaurus.[4] A year earlier, Lawrence Lambe described the short, two-fingered forelimbs of the closely related Gorgosaurus.[110] This strongly suggested that T. rex had similar forelimbs, but this hypothesis was not confirmed until the first complete T. rex forelimbs were identified in 1989, belonging to MOR 555 (the "Wankel rex").[111][112] The remains of Sue also include complete forelimbs.[47] T. rex arms are very small relative to overall body size, measuring only 1 meter (3.3 ft) long, and some scholars have labelled them as vestigial. The bones show large areas for muscle attachment, indicating considerable strength. This was recognized as early as 1906 by Osborn, who speculated that the forelimbs may have been used to grasp a mate during copulation.[8] It has also been suggested that the forelimbs were used to assist the animal in rising from a prone position.[104]
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+ Another possibility is that the forelimbs held struggling prey while it was killed by the tyrannosaur's enormous jaws. This hypothesis may be supported by biomechanical analysis. T. rex forelimb bones exhibit extremely thick cortical bone, which has been interpreted as evidence that they were developed to withstand heavy loads. The biceps brachii muscle of an adult T. rex was capable of lifting 199 kilograms (439 lb) by itself; other muscles such as the brachialis would work along with the biceps to make elbow flexion even more powerful. The M. biceps muscle of T. rex was 3.5 times as powerful as the human equivalent. A T. rex forearm had a limited range of motion, with the shoulder and elbow joints allowing only 40 and 45 degrees of motion, respectively. In contrast, the same two joints in Deinonychus allow up to 88 and 130 degrees of motion, respectively, while a human arm can rotate 360 degrees at the shoulder and move through 165 degrees at the elbow. The heavy build of the arm bones, strength of the muscles, and limited range of motion may indicate a system evolved to hold fast despite the stresses of a struggling prey animal. In the first detailed scientific description of Tyrannosaurus forelimbs, paleontologists Kenneth Carpenter and Matt Smith dismissed notions that the forelimbs were useless or that T. rex was an obligate scavenger.[113]
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+ According to paleontologist Steven M. Stanley, the 1 metre (3.3 ft) arms of T. rex were used for slashing prey, especially by using its claws to rapidly inflict long, deep gashes to its prey, although this concept is disputed by others believing the arms were used for grasping a sexual partner.[114]
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+ As of 2014, it is not clear if Tyrannosaurus was endothermic (“warm-blooded”). Tyrannosaurus, like most dinosaurs, was long thought to have an ectothermic ("cold-blooded") reptilian metabolism. The idea of dinosaur ectothermy was challenged by scientists like Robert T. Bakker and John Ostrom in the early years of the "Dinosaur Renaissance", beginning in the late 1960s.[115][116] T. rex itself was claimed to have been endothermic ("warm-blooded"), implying a very active lifestyle.[35] Since then, several paleontologists have sought to determine the ability of Tyrannosaurus to regulate its body temperature. Histological evidence of high growth rates in young T. rex, comparable to those of mammals and birds, may support the hypothesis of a high metabolism. Growth curves indicate that, as in mammals and birds, T. rex growth was limited mostly to immature animals, rather than the indeterminate growth seen in most other vertebrates.[89]
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+ Oxygen isotope ratios in fossilized bone are sometimes used to determine the temperature at which the bone was deposited, as the ratio between certain isotopes correlates with temperature. In one specimen, the isotope ratios in bones from different parts of the body indicated a temperature difference of no more than 4 to 5 °C (7 to 9 °F) between the vertebrae of the torso and the tibia of the lower leg. This small temperature range between the body core and the extremities was claimed by paleontologist Reese Barrick and geochemist William Showers to indicate that T. rex maintained a constant internal body temperature (homeothermy) and that it enjoyed a metabolism somewhere between ectothermic reptiles and endothermic mammals.[117] Other scientists have pointed out that the ratio of oxygen isotopes in the fossils today does not necessarily represent the same ratio in the distant past, and may have been altered during or after fossilization (diagenesis).[118] Barrick and Showers have defended their conclusions in subsequent papers, finding similar results in another theropod dinosaur from a different continent and tens of millions of years earlier in time (Giganotosaurus).[119] Ornithischian dinosaurs also showed evidence of homeothermy, while varanid lizards from the same formation did not.[120] Even if T. rex does exhibit evidence of homeothermy, it does not necessarily mean that it was endothermic. Such thermoregulation may also be explained by gigantothermy, as in some living sea turtles.[121][122][123] Similar to contemporary alligators, dorsotemporal fenestra in Tyrannosaurus's skull may have aided thermoregulation.[124]
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+ In the March 2005 issue of Science, Mary Higby Schweitzer of North Carolina State University and colleagues announced the recovery of soft tissue from the marrow cavity of a fossilized leg bone from a T. rex. The bone had been intentionally, though reluctantly, broken for shipping and then not preserved in the normal manner, specifically because Schweitzer was hoping to test it for soft tissue.[125] Designated as the Museum of the Rockies specimen 1125, or MOR 1125, the dinosaur was previously excavated from the Hell Creek Formation. Flexible, bifurcating blood vessels and fibrous but elastic bone matrix tissue were recognized. In addition, microstructures resembling blood cells were found inside the matrix and vessels. The structures bear resemblance to ostrich blood cells and vessels. Whether an unknown process, distinct from normal fossilization, preserved the material, or the material is original, the researchers do not know, and they are careful not to make any claims about preservation.[126] If it is found to be original material, any surviving proteins may be used as a means of indirectly guessing some of the DNA content of the dinosaurs involved, because each protein is typically created by a specific gene. The absence of previous finds may be the result of people assuming preserved tissue was impossible, therefore not looking. Since the first, two more tyrannosaurs and a hadrosaur have also been found to have such tissue-like structures.[125] Research on some of the tissues involved has suggested that birds are closer relatives to tyrannosaurs than other modern animals.[127]
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+ In studies reported in Science in April 2007, Asara and colleagues concluded that seven traces of collagen proteins detected in purified T. rex bone most closely match those reported in chickens, followed by frogs and newts. The discovery of proteins from a creature tens of millions of years old, along with similar traces the team found in a mastodon bone at least 160,000 years old, upends the conventional view of fossils and may shift paleontologists' focus from bone hunting to biochemistry. Until these finds, most scientists presumed that fossilization replaced all living tissue with inert minerals. Paleontologist Hans Larsson of McGill University in Montreal, who was not part of the studies, called the finds "a milestone", and suggested that dinosaurs could "enter the field of molecular biology and really slingshot paleontology into the modern world".[128]
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+ The presumed soft tissue was called into question by Thomas Kaye of the University of Washington and his co-authors in 2008. They contend that what was really inside the tyrannosaur bone was slimy biofilm created by bacteria that coated the voids once occupied by blood vessels and cells.[129] The researchers found that what previously had been identified as remnants of blood cells, because of the presence of iron, were actually framboids, microscopic mineral spheres bearing iron. They found similar spheres in a variety of other fossils from various periods, including an ammonite. In the ammonite they found the spheres in a place where the iron they contain could not have had any relationship to the presence of blood.[130] Schweitzer has strongly criticized Kaye's claims and argues that there is no reported evidence that biofilms can produce branching, hollow tubes like those noted in her study.[131] San Antonio, Schweitzer and colleagues published an analysis in 2011 of what parts of the collagen had been recovered, finding that it was the inner parts of the collagen coil that had been preserved, as would have been expected from a long period of protein degradation.[132] Other research challenges the identification of soft tissue as biofilm and confirms finding "branching, vessel-like structures" from within fossilized bone.[133]
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+ Scientists have produced a wide range of possible maximum running speeds for Tyrannosaurus, mostly around 11 meters per second (40 km/h; 25 mph), but as low as 5–11 meters per second (18–40 km/h; 11–25 mph) and as high as 20 meters per second (72 km/h; 45 mph). Estimates that Tyrannosaurus had relatively larger leg muscles than any animal alive today indicate that fast running was possible at speeds of 40–70 kilometers per hour (25–43 mph).[134] Researchers have relied on various estimating techniques because, while there are many tracks of large theropods walking, none had the pattern of running.[135]
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+ A 2002 report used a mathematical model (validated by applying it to three living animals: alligators, chickens, and humans; and eight more species, including emus and ostriches[135]) to gauge the leg muscle mass needed for fast running (over 40 km/h or 25 mph).[134] Scientists who think that Tyrannosaurus was able to run point out that hollow bones and other features that would have lightened its body may have kept adult weight to a mere 4.5 metric tons (5.0 short tons) or so, or that other animals like ostriches and horses with long, flexible legs are able to achieve high speeds through slower but longer strides.[135] Proposed top speeds exceeded 40 kilometers per hour (25 mph) for Tyrannosaurus, but were deemed infeasible because they would require exceptional leg muscles of approximately 40–86% of total body mass. Even moderately fast speeds would have required large leg muscles. If the muscle mass was less, only 18 kilometers per hour (11 mph) for walking or jogging would have been possible.[134] Holtz noted that tyrannosaurids and some closely related groups had significantly longer distal hindlimb components (shin plus foot plus toes) relative to the femur length than most other theropods, and that tyrannosaurids and their close relatives had a tightly interlocked metatarsus (foot bones).[136] The third metatarsal was squeezed between the second and fourth metatarsals to form a single unit called an arctometatarsus. This ankle feature may have helped the animal to run more efficiently.[137] Together, these leg features allowed Tyrannosaurus to transmit locomotory forces from the foot to the lower leg more effectively than in earlier theropods.[136]
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+ Additionally, a 2020 study indicates that Tyrannosaurus and other tyrannosaurids were exceptionally efficient walkers. Studies by Dececchi et al., compared the leg proportions, body mass, and the gaits of more than 70 species of theropod dinosaurs including Tyrannosaurus and its relatives. The research team then applied a variety of methods to estimate each dinosaur's top speed when running as well as how much energy each dinosaur expended while moving at more relaxed speeds such as when walking. Among smaller to medium-sized species such as dromaeosaurids, longer legs appear to be an adaptation for faster running, in line with previous results by other researchers. But for theropods weighing over 1,000 kg (2,200 lb), top running speed is limited by body size, so longer legs instead were found to have correlated with low-energy walking. The results further indicate that smaller theropods evolved long legs as a means to both aid in hunting and escape from larger predators while larger theropods that evolved long legs did so to reduce the energy costs and increase foraging efficiency, as they were freed from the demands of predation pressure due to their role as apex predators. Compared to more basal groups of theropods in the study, tyrannosaurs like Tyrannosaurus itself showed a marked increase in foraging efficiency due to reduced energy expenditures during hunting or scavenging. This in turn likely resulted in tyrannosaurs having a reduced need for hunting forays and requiring less food to sustain themselves as a result. Additionally, the research, in conjunction with studies that show tyrannosaurs were more agile than other large bodied-theropods, indicates they were quite well-adapted to a long-distance stalking approach followed by a quick burst of speed to go for the kill. Analogies can be noted between tyrannosaurids and modern wolves as a result, supported by evidence that at least some tyrannosaurids were hunting in group settings.[138][139]
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+ A 2017 study estimated the top running speed of Tyrannosaurus as 17 mph (27 km/h), speculating that Tyrannosaurus exhausted its energy reserves long before reaching top speed, resulting in a parabola-like relationship between size and speed.[140][141] Another 2017 study hypothesized that an adult Tyrannosaurus was incapable of running due to high skeletal loads. Using a calculated weight estimate of 7 tons, the model showed that speeds above 11 mph (18 km/h) would have probably shattered the leg bones of Tyrannosaurus. The finding may mean that running was also not possible for other giant theropod dinosaurs like Giganotosaurus, Mapusaurus and Acrocanthosaurus.[142]
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+ However, studies by Eric Snively and colleagues, published in 2019 indicate that Tyrannosaurus and other tyrannosaurids were more maneuverable than allosauroids and other theropods of comparable size due to low rotational inertia compared to their body mass combined with large leg muscles. As a result, it is hypothesized that Tyrannosaurus was capable of making relatively quick turns and could likely pivot its body more quickly when close to its prey, or that while turning, the theropod could "pirouette" on a single planted foot while the alternating leg was held out in a suspended swing during pursuit. The results of this study potentially could shed light on how agility could have contributed to the success of tyrannosaurid evolution.[143]
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+ A study conducted by Lawrence Witmer and Ryan Ridgely of Ohio University found that Tyrannosaurus shared the heightened sensory abilities of other coelurosaurs, highlighting relatively rapid and coordinated eye and head movements; an enhanced ability to sense low frequency sounds, which would allow tyrannosaurs to track prey movements from long distances; and an enhanced sense of smell.[144] A study published by Kent Stevens concluded that Tyrannosaurus had keen vision. By applying modified perimetry to facial reconstructions of several dinosaurs including Tyrannosaurus, the study found that Tyrannosaurus had a binocular range of 55 degrees, surpassing that of modern hawks. Stevens estimated that Tyrannosaurus had 13 times the visual acuity of a human and surpassed the visual acuity of an eagle, which is 3.6 times that of a person. Stevens estimated a limiting far point (that is, the distance at which an object can be seen as separate from the horizon) as far as 6 km (3.7 mi) away, which is greater than the 1.6 km (1 mi) that a human can see.[41][42][145]
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+ Thomas Holtz Jr. would note that high depth perception of Tyrannosaurus may have been due to the prey it had to hunt, noting that it had to hunt horned dinosaurs such as Triceratops, armored dinosaurs such as Ankylosaurus, and the duck-billed dinosaurs and their possibly complex social behaviors. He would suggest that this made precision more crucial for Tyrannosaurus enabling it to, "get in, get that blow in and take it down." In contrast, Acrocanthosaurus had limited depth perception because they hunted large sauropods, which were relatively rare during the time of Tyrannosaurus.[96]
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+ Tyrannosaurus had very large olfactory bulbs and olfactory nerves relative to their brain size, the organs responsible for a heightened sense of smell. This suggests that the sense of smell was highly developed, and implies that tyrannosaurs could detect carcasses by scent alone across great distances. The sense of smell in tyrannosaurs may have been comparable to modern vultures, which use scent to track carcasses for scavenging. Research on the olfactory bulbs has shown that T. rex had the most highly developed sense of smell of 21 sampled non-avian dinosaur species.[146]
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+
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+ Somewhat unusually among theropods, T. rex had a very long cochlea. The length of the cochlea is often related to hearing acuity, or at least the importance of hearing in behavior, implying that hearing was a particularly important sense to tyrannosaurs. Specifically, data suggests that T. rex heard best in the low-frequency range, and that low-frequency sounds were an important part of tyrannosaur behavior.[144] A 2017 study by Thomas Carr and colleagues found that the snout of tyrannosaurids was highly sensitive, based on a high number of small openings in the facial bones of the related Daspletosaurus that contained sensory neurons. The study speculated that tyrannosaurs might have used their sensitive snouts to measure the temperature of their nests and to gently pick-up eggs and hatchlings, as seen in modern crocodylians.[51]
157
+
158
+ A study by Grant R. Hurlburt, Ryan C. Ridgely and Lawrence Witmer obtained estimates for Encephalization Quotients (EQs), based on reptiles and birds, as well as estimates for the ratio of cerebrum to brain mass. The study concluded that Tyrannosaurus had the relatively largest brain of all adult non-avian dinosaurs with the exception of certain small maniraptoriforms (Bambiraptor, Troodon and Ornithomimus). The study found that Tyrannosaurus's relative brain size was still within the range of modern reptiles, being at most 2 standard deviations above the mean of non-avian reptile EQs. The estimates for the ratio of cerebrum mass to brain mass would range from 47.5 to 49.53 percent. According to the study, this is more than the lowest estimates for extant birds (44.6 percent), but still close to the typical ratios of the smallest sexually mature alligators which range from 45.9–47.9 percent.[147] Other studies, such as those by Steve Brusatte, indicate the encephalization quotient of Tyrannosaurus was similar in range (2.0-2.4) to a chimpanzee (2.2-2.5), though this may be debatable as reptilian and mammalian encephalization quotients are not equivalent.[148]
159
+
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+ Suggesting that Tyrannosaurus may have been pack hunters, Philip J. Currie compared T. rex to related species Tarbosaurus bataar and Albertosaurus sarcophagus, citing fossil evidence that may indicate pack behavior.[149] A find in South Dakota where three T. rex skeletons were in close proximity suggested a pack.[150][151] Because available prey such as Triceratops and Ankylosaurus had significant defenses, it may have been effective for T. rex to hunt in groups.[149]
161
+
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+ Currie's pack-hunting hypothesis has been criticized for not having been peer-reviewed, but rather was discussed in a television interview and book called Dino Gangs.[152] The Currie theory for pack hunting by T. rex is based mainly by analogy to a different species, Tarbosaurus bataar, and that the supposed evidence for pack hunting in T. bataar itself had not yet been peer-reviewed. According to scientists assessing the Dino Gangs program, the evidence for pack hunting in Tarbosaurus and Albertosaurus is weak and based on skeletal remains for which alternate explanations may apply (such as drought or a flood forcing dinosaurs to die together in one place).[152] Fossilized trackways from the Upper Cretaceous Wapiti Formation of northeastern British Columbia, Canada, left by three tyrannosaurids traveling in the same direction, may also indicate packs.[153][154]
163
+
164
+ Evidence of intraspecific attack were found by Joseph Peterson and his colleagues in the juvenile Tyrannosaurus nicknamed Jane. Peterson and his team found that Jane's skull showed healed puncture wounds on the upper jaw and snout which they believe came from another juvenile Tyrannosaurus. Subsequent CT scans of Jane's skull would further confirm the team's hypothesis, showing that the puncture wounds came from a traumatic injury and that there was subsequent healing.[155] The team would also state that Jane's injuries were structurally different from the parasite-induced lesions found in Sue and that Jane's injuries were on her face whereas the parasite that infected Sue caused lesions to the lower jaw.[156]
165
+
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+ Most paleontologists accept that Tyrannosaurus was both an active predator and a scavenger like most large carnivores.[157] By far the largest carnivore in its environment, T. rex was most likely an apex predator, preying upon hadrosaurs, armored herbivores like ceratopsians and ankylosaurs, and possibly sauropods.[158] A study in 2012 by Karl Bates and Peter Falkingham found that Tyrannosaurus had the most powerful bite of any terrestrial animal that has ever lived, finding an adult Tyrannosaurus could have exerted 35,000 to 57,000 N (7,868 to 12,814 lbf) of force in the back teeth.[159][160][161] Even higher estimates were made by Mason B. Meers in 2003.[44] This allowed it to crush bones during repetitive biting and fully consume the carcasses of large dinosaurs.[20] Stephan Lautenschlager and colleagues calculated that Tyrannosaurus was capable of a maximum jaw gape of around 80 degrees, a necessary adaptation for a wide range of jaw angles to power the creature's strong bite.[162][163]
167
+
168
+ A debate exists, however, about whether Tyrannosaurus was primarily a predator or a pure scavenger; the debate was assessed in a 1917 study by Lambe which argued Tyrannosaurus was a pure scavenger because the Gorgosaurus teeth showed hardly any wear.[164] This argument may not be valid because theropods replaced their teeth quite rapidly. Ever since the first discovery of Tyrannosaurus most scientists have speculated that it was a predator; like modern large predators it would readily scavenge or steal another predator's kill if it had the opportunity.[165]
169
+
170
+ Paleontologist Jack Horner has been a major proponent of the view that Tyrannosaurus was not a predator at all but instead was exclusively a scavenger.[111][166][167] He has put forward arguments in the popular literature to support the pure scavenger hypothesis:
171
+
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+ Other evidence suggests hunting behavior in Tyrannosaurus. The eye sockets of tyrannosaurs are positioned so that the eyes would point forward, giving them binocular vision slightly better than that of modern hawks. It is not obvious why natural selection would have favored this long-term trend if tyrannosaurs had been pure scavengers, which would not have needed the advanced depth perception that stereoscopic vision provides.[41][42] In modern animals, binocular vision is found mainly in predators.
173
+
174
+ A skeleton of the hadrosaurid Edmontosaurus annectens has been described from Montana with healed tyrannosaur-inflicted damage on its tail vertebrae. The fact that the damage seems to have healed suggests that the Edmontosaurus survived a tyrannosaur's attack on a living target, i.e. the tyrannosaur had attempted active predation.[169] Despite the consensus that the tail bites were caused by Tyrannosaurus, there has been some evidence to show that they might have been created by other factors. For example, a 2014 study suggested that the tail injuries might have been due to Edmontosaurus indivduals stepping on each other,[170] while another study in 2020 backs up the hypothesis that biomechanical stress is the cause for the tail injuries.[171]. There is also evidence for an aggressive interaction between a Triceratops and a Tyrannosaurus in the form of partially healed tyrannosaur tooth marks on a Triceratops brow horn and squamosal (a bone of the neck frill); the bitten horn is also broken, with new bone growth after the break. It is not known what the exact nature of the interaction was, though: either animal could have been the aggressor.[172] Since the Triceratops wounds healed, it is most likely that the Triceratops survived the encounter and managed to overcome the Tyrannosaurus. In a battle against a bull Triceratops, the Triceratops would likely defend itself by inflicting fatal wounds to the Tyrannosaurus using its sharp horns.[173] Studies of Sue found a broken and healed fibula and tail vertebrae, scarred facial bones and a tooth from another Tyrannosaurus embedded in a neck vertebra, providing evidence for aggressive behavior.[174] Studies on hadrosaur vertebrae from the Hell Creek Formation that were punctured by the teeth of what appears to be a late-stage juvenile Tyrannosaurus indicate that despite lacking the bone-crushing adaptations of the adults, young individuals were still capable of using the same bone-puncturing feeding technique as their adult counterparts.[175]
175
+
176
+ Tyrannosaurus may have had infectious saliva used to kill its prey, as proposed by William Abler in 1992. Abler observed that the serrations (tiny protuberances) on the cutting edges of the teeth are closely spaced, enclosing little chambers. These chambers might have trapped pieces of carcass with bacteria, giving Tyrannosaurus a deadly, infectious bite much like the Komodo dragon was thought to have.[176][177] Jack Horner and Don Lessem, in a 1993 popular book, questioned Abler's hypothesis, arguing that Tyrannosaurus's tooth serrations as more like cubes in shape than the serrations on a Komodo monitor's teeth, which are rounded.[111]:214–215
177
+
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+ Tyrannosaurus, and most other theropods, probably primarily processed carcasses with lateral shakes of the head, like crocodilians. The head was not as maneuverable as the skulls of allosauroids, due to flat joints of the neck vertebrae.[178]
179
+
180
+ In 2001, Bruce Rothschild and others published a study examining evidence for stress fractures and tendon avulsions in theropod dinosaurs and the implications for their behavior. Since stress fractures are caused by repeated trauma rather than singular events they are more likely to be caused by regular behavior than other types of injuries. Of the 81 Tyrannosaurus foot bones examined in the study one was found to have a stress fracture, while none of the 10 hand bones were found to have stress fractures. The researchers found tendon avulsions only among Tyrannosaurus and Allosaurus. An avulsion injury left a divot on the humerus of Sue the T. rex, apparently located at the origin of the deltoid or teres major muscles. The presence of avulsion injuries being limited to the forelimb and shoulder in both Tyrannosaurus and Allosaurus suggests that theropods may have had a musculature more complex than and functionally different from those of birds. The researchers concluded that Sue's tendon avulsion was probably obtained from struggling prey. The presence of stress fractures and tendon avulsions in general provides evidence for a "very active" predation-based diet rather than obligate scavenging.[179]
181
+
182
+ A 2009 study showed that smooth-edged holes in the skulls of several specimens might have been caused by Trichomonas-like parasites that commonly infect birds. Seriously infected individuals, including "Sue" and MOR 980 ("Peck's Rex"), might therefore have died from starvation after feeding became increasingly difficult. Previously, these holes had been explained by the bacterious bone infection Actinomycosis or by intraspecific attacks.[180]
183
+
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+ One study of Tyrannosaurus specimens with tooth marks in the bones attributable to the same genus was presented as evidence of cannibalism.[181] Tooth marks in the humerus, foot bones and metatarsals, may indicate opportunistic scavenging, rather than wounds caused by combat with another T. rex.[181][182] Other tyrannosaurids may also have practiced cannibalism.[181]
185
+
186
+ Tyrannosaurus lived during what is referred to as the Lancian faunal stage (Maastrichtian age) at the end of the Late Cretaceous. Tyrannosaurus ranged from Canada in the north to at least New Mexico in the south of Laramidia.[5] During this time Triceratops was the major herbivore in the northern portion of its range, while the titanosaurian sauropod Alamosaurus "dominated" its southern range. Tyrannosaurus remains have been discovered in different ecosystems, including inland and coastal subtropical, and semi-arid plains.
187
+
188
+ Several notable Tyrannosaurus remains have been found in the Hell Creek Formation. During the Maastrichtian this area was subtropical, with a warm and humid climate. The flora consisted mostly of angiosperms, but also included trees like dawn redwood (Metasequoia) and Araucaria. Tyrannosaurus shared this ecosystem with ceratopsians Leptoceratops, Torosaurus, and Triceratops, the hadrosaurid Edmontosaurus annectens, the parksosaurid Thescelosaurus, the ankylosaurs Ankylosaurus and Denversaurus, the pachycephalosaurs Pachycephalosaurus and Sphaerotholus, and the theropods Ornithomimus, Struthiomimus, Acheroraptor, Dakotaraptor, Pectinodon and Anzu.[183]
189
+
190
+ Another formation with Tyrannosaurus remains is the Lance Formation of Wyoming. This has been interpreted as a bayou environment similar to today's Gulf Coast. The fauna was very similar to Hell Creek, but with Struthiomimus replacing its relative Ornithomimus. The small ceratopsian Leptoceratops also lived in the area.[184]
191
+
192
+ In its southern range Tyrannosaurus lived alongside the titanosaur Alamosaurus, the ceratopsians Torosaurus, Bravoceratops and Ojoceratops, hadrosaurs which consisted of a species of Edmontosaurus, Kritosaurus and a possible species of Gryposaurus, the nodosaur Glyptodontopelta, the oviraptorid Ojoraptosaurus, possible species of the theropods Troodon and Richardoestesia, and the pterosaur Quetzalcoatlus.[185] The region is thought to have been dominated by semi-arid inland plains, following the probable retreat of the Western Interior Seaway as global sea levels fell.[186]
193
+
194
+ Tyrannosaurus may have also inhabited Mexico's Lomas Coloradas formation in Sonora. Though skeletal evidence is lacking, six shed and broken teeth from the fossil bed have been thoroughly compared with other theropod genera and appear to be identical to those of Tyrannosaurus. If true, the evidence indicates the range of Tyrannosaurus was possibly more extensive than previously believed.[187] It is possible that tyrannosaurs were originally Asian species, migrating to North America before the end of the Cretaceous period.[188]
195
+
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+ Since it was first described in 1905, T. rex has become the most widely recognized dinosaur species in popular culture. It is the only dinosaur that is commonly known to the general public by its full scientific name (binomial name) and the scientific abbreviation T. rex has also come into wide usage.[47] Robert T. Bakker notes this in The Dinosaur Heresies and explains that, "a name like ‘T. rex’ is just irresistible to the tongue."[35]
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en/5786.html.txt ADDED
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1
+
2
+
3
+ A triangle is a polygon with three edges and three vertices. It is one of the basic shapes in geometry. A triangle with vertices A, B, and C is denoted
4
+
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+
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+
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+ △
8
+ A
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+ B
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+ C
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+
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+
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+ {\displaystyle \triangle ABC}
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+
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+ .
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+
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+ In Euclidean geometry any three points, when non-collinear, determine a unique triangle and simultaneously, a unique plane (i.e. a two-dimensional Euclidean space). In other words, there is only one plane that contains that triangle, and every triangle is contained in some plane. If the entire geometry is only the Euclidean plane, there is only one plane and all triangles are contained in it; however, in higher-dimensional Euclidean spaces, this is no longer true. This article is about triangles in Euclidean geometry, and in particular, the Euclidean plane, except where otherwise noted.
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+ Triangles can be classified according to the lengths of their sides:
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+
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+ Hatch marks, also called tick marks, are used in diagrams of triangles and other geometric figures to identify sides of equal lengths. A side can be marked with a pattern of "ticks", short line segments in the form of tally marks; two sides have equal lengths if they are both marked with the same pattern. In a triangle, the pattern is usually no more than 3 ticks. An equilateral triangle has the same pattern on all 3 sides, an isosceles triangle has the same pattern on just 2 sides, and a scalene triangle has different patterns on all sides since no sides are equal. Similarly, patterns of 1, 2, or 3 concentric arcs inside the angles are used to indicate equal angles. An equilateral triangle has the same pattern on all 3 angles, an isosceles triangle has the same pattern on just 2 angles, and a scalene triangle has different patterns on all angles since no angles are equal.
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+
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+ Triangles can also be classified according to their internal angles, measured here in degrees.
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+
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+ A triangle that has two angles with the same measure also has two sides with the same length, and therefore it is an isosceles triangle. It follows that in a triangle where all angles have the same measure, all three sides have the same length, and such a triangle is therefore equilateral.
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+ Triangles are assumed to be two-dimensional plane figures, unless the context provides otherwise (see Non-planar triangles, below). In rigorous treatments, a triangle is therefore called a 2-simplex (see also Polytope). Elementary facts about triangles were presented by Euclid in books 1–4 of his Elements, around 300 BC.
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+
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+ The sum of the measures of the interior angles of a triangle in Euclidean space is always 180 degrees.[5] This fact is equivalent to Euclid's parallel postulate. This allows determination of the measure of the third angle of any triangle given the measure of two angles. An exterior angle of a triangle is an angle that is a linear pair (and hence supplementary) to an interior angle. The measure of an exterior angle of a triangle is equal to the sum of the measures of the two interior angles that are not adjacent to it; this is the exterior angle theorem. The sum of the measures of the three exterior angles (one for each vertex) of any triangle is 360 degrees.[note 2]
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+
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+ Two triangles are said to be similar if every angle of one triangle has the same measure as the corresponding angle in the other triangle. The corresponding sides of similar triangles have lengths that are in the same proportion, and this property is also sufficient to establish similarity.
32
+
33
+ Some basic theorems about similar triangles are:
34
+
35
+ Two triangles that are congruent have exactly the same size and shape:[note 4] all pairs of corresponding interior angles are equal in measure, and all pairs of corresponding sides have the same length. (This is a total of six equalities, but three are often sufficient to prove congruence.)
36
+
37
+ Some individually necessary and sufficient conditions for a pair of triangles to be congruent are:
38
+
39
+ Some individually sufficient conditions are:
40
+
41
+ An important condition is:
42
+
43
+ Using right triangles and the concept of similarity, the trigonometric functions sine and cosine can be defined. These are functions of an angle which are investigated in trigonometry.
44
+
45
+ A central theorem is the Pythagorean theorem, which states in any right triangle, the square of the length of the hypotenuse equals the sum of the squares of the lengths of the two other sides. If the hypotenuse has length c, and the legs have lengths a and b, then the theorem states that
46
+
47
+ The converse is true: if the lengths of the sides of a triangle satisfy the above equation, then the triangle has a right angle opposite side c.
48
+
49
+ Some other facts about right triangles:
50
+
51
+ For all triangles, angles and sides are related by the law of cosines and law of sines (also called the cosine rule and sine rule).
52
+
53
+ The triangle inequality states that the sum of the lengths of any two sides of a triangle must be greater than or equal to the length of the third side. That sum can equal the length of the third side only in the case of a degenerate triangle, one with collinear vertices. It is not possible for that sum to be less than the length of the third side. A triangle with three given positive side lengths exists if and only if those side lengths satisfy the triangle inequality.
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+
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+ Three given angles form a non-degenerate triangle (and indeed an infinitude of them) if and only if both of these conditions hold: (a) each of the angles is positive, and (b) the angles sum to 180°. If degenerate triangles are permitted, angles of 0° are permitted.
56
+
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+ Three positive angles α, β, and γ, each of them less than 180°, are the angles of a triangle if and only if any one of the following conditions holds:
58
+
59
+ the last equality applying only if none of the angles is 90° (so the tangent function's value is always finite).
60
+
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+ There are thousands of different constructions that find a special point associated with (and often inside) a triangle, satisfying some unique property: see the article Encyclopedia of Triangle Centers for a catalogue of them. Often they are constructed by finding three lines associated in a symmetrical way with the three sides (or vertices) and then proving that the three lines meet in a single point: an important tool for proving the existence of these is Ceva's theorem, which gives a criterion for determining when three such lines are concurrent. Similarly, lines associated with a triangle are often constructed by proving that three symmetrically constructed points are collinear: here Menelaus' theorem gives a useful general criterion. In this section just a few of the most commonly encountered constructions are explained.
62
+
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+ A perpendicular bisector of a side of a triangle is a straight line passing through the midpoint of the side and being perpendicular to it, i.e. forming a right angle with it. The three perpendicular bisectors meet in a single point, the triangle's circumcenter, usually denoted by O; this point is the center of the circumcircle, the circle passing through all three vertices. The diameter of this circle, called the circumdiameter, can be found from the law of sines stated above. The circumcircle's radius is called the circumradius.
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+
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+ Thales' theorem implies that if the circumcenter is located on a side of the triangle, then the opposite angle is a right one. If the circumcenter is located inside the triangle, then the triangle is acute; if the circumcenter is located outside the triangle, then the triangle is obtuse.
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+
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+ An altitude of a triangle is a straight line through a vertex and perpendicular to (i.e. forming a right angle with) the opposite side. This opposite side is called the base of the altitude, and the point where the altitude intersects the base (or its extension) is called the foot of the altitude. The length of the altitude is the distance between the base and the vertex. The three altitudes intersect in a single point, called the orthocenter of the triangle, usually denoted by H. The orthocenter lies inside the triangle if and only if the triangle is acute.
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+
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+ An angle bisector of a triangle is a straight line through a vertex which cuts the corresponding angle in half. The three angle bisectors intersect in a single point, the incenter, usually denoted by I, the center of the triangle's incircle. The incircle is the circle which lies inside the triangle and touches all three sides. Its radius is called the inradius. There are three other important circles, the excircles; they lie outside the triangle and touch one side as well as the extensions of the other two. The centers of the in- and excircles form an orthocentric system.
70
+
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+ A median of a triangle is a straight line through a vertex and the midpoint of the opposite side, and divides the triangle into two equal areas. The three medians intersect in a single point, the triangle's centroid or geometric barycenter, usually denoted by G. The centroid of a rigid triangular object (cut out of a thin sheet of uniform density) is also its center of mass: the object can be balanced on its centroid in a uniform gravitational field. The centroid cuts every median in the ratio 2:1, i.e. the distance between a vertex and the centroid is twice the distance between the centroid and the midpoint of the opposite side.
72
+
73
+ The midpoints of the three sides and the feet of the three altitudes all lie on a single circle, the triangle's nine-point circle. The remaining three points for which it is named are the midpoints of the portion of altitude between the vertices and the orthocenter. The radius of the nine-point circle is half that of the circumcircle. It touches the incircle (at the Feuerbach point) and the three excircles.
74
+
75
+ The orthocenter (blue point), center of the nine-point circle (red), centroid (orange), and circumcenter (green) all lie on a single line, known as Euler's line (red line). The center of the nine-point circle lies at the midpoint between the orthocenter and the circumcenter, and the distance between the centroid and the circumcenter is half that between the centroid and the orthocenter.
76
+
77
+ The center of the incircle is not in general located on Euler's line.
78
+
79
+ If one reflects a median in the angle bisector that passes through the same vertex, one obtains a symmedian. The three symmedians intersect in a single point, the symmedian point of the triangle.
80
+
81
+ There are various standard methods for calculating the length of a side or the measure of an angle. Certain methods are suited to calculating values in a right-angled triangle; more complex methods may be required in other situations.
82
+
83
+ In right triangles, the trigonometric ratios of sine, cosine and tangent can be used to find unknown angles and the lengths of unknown sides. The sides of the triangle are known as follows:
84
+
85
+ The sine of an angle is the ratio of the length of the opposite side to the length of the hypotenuse. In our case
86
+
87
+ This ratio does not depend on the particular right triangle chosen, as long as it contains the angle A, since all those triangles are similar.
88
+
89
+ The cosine of an angle is the ratio of the length of the adjacent side to the length of the hypotenuse. In our case
90
+
91
+ The tangent of an angle is the ratio of the length of the opposite side to the length of the adjacent side. In our case
92
+
93
+ The acronym "SOH-CAH-TOA" is a useful mnemonic for these ratios.
94
+
95
+ The inverse trigonometric functions can be used to calculate the internal angles for a right angled triangle with the length of any two sides.
96
+
97
+ Arcsin can be used to calculate an angle from the length of the opposite side and the length of the hypotenuse.
98
+
99
+ Arccos can be used to calculate an angle from the length of the adjacent side and the length of the hypotenuse.
100
+
101
+ Arctan can be used to calculate an angle from the length of the opposite side and the length of the adjacent side.
102
+
103
+ In introductory geometry and trigonometry courses, the notation sin−1, cos−1, etc., are often used in place of arcsin, arccos, etc. However, the arcsin, arccos, etc., notation is standard in higher mathematics where trigonometric functions are commonly raised to powers, as this avoids confusion between multiplicative inverse and compositional inverse.
104
+
105
+ The law of sines, or sine rule,[8] states that the ratio of the length of a side to the sine of its corresponding opposite angle is constant, that is
106
+
107
+ This ratio is equal to the diameter of the circumscribed circle of the given triangle. Another interpretation of this theorem is that every triangle with angles α, β and γ is similar to a triangle with side lengths equal to sin α, sin β and sin γ. This triangle can be constructed by first constructing a circle of diameter 1, and inscribing in it two of the angles of the triangle. The length of the sides of that triangle will be sin α, sin β and sin γ. The side whose length is sin α is opposite to the angle whose measure is α, etc.
108
+
109
+ The law of cosines, or cosine rule, connects the length of an unknown side of a triangle to the length of the other sides and the angle opposite to the unknown side.[8] As per the law:
110
+
111
+ For a triangle with length of sides a, b, c and angles of α, β, γ respectively, given two known lengths of a triangle a and b, and the angle between the two known sides γ (or the angle opposite to the unknown side c), to calculate the third side c, the following formula can be used:
112
+
113
+ If the lengths of all three sides of any triangle are known the three angles can be calculated:
114
+
115
+ The law of tangents, or tangent rule, can be used to find a side or an angle when two sides and an angle or two angles and a side are known. It states that:[9]
116
+
117
+ "Solution of triangles" is the main trigonometric problem: to find missing characteristics of a triangle (three angles, the lengths of the three sides etc.) when at least three of these characteristics are given. The triangle can be located on a plane or on a sphere. This problem often occurs in various trigonometric applications, such as geodesy, astronomy, construction, navigation etc.
118
+
119
+ Calculating the area T of a triangle is an elementary problem encountered often in many different situations. The best known and simplest formula is:
120
+
121
+ where b is the length of the base of the triangle, and h is the height or altitude of the triangle. The term "base" denotes any side, and "height" denotes the length of a perpendicular from the vertex opposite the base onto the line containing the base. In 499 CE Aryabhata, used this illustrated method in the Aryabhatiya (section 2.6).[10]
122
+
123
+ Although simple, this formula is only useful if the height can be readily found, which is not always the case. For example, the surveyor of a triangular field might find it relatively easy to measure the length of each side, but relatively difficult to construct a 'height'. Various methods may be used in practice, depending on what is known about the triangle. The following is a selection of frequently used formulae for the area of a triangle.[11]
124
+
125
+ The height of a triangle can be found through the application of trigonometry.
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+
127
+ Knowing SAS: Using the labels in the image on the right, the altitude is h = a sin
128
+
129
+
130
+
131
+ γ
132
+
133
+
134
+ {\displaystyle \gamma }
135
+
136
+ . Substituting this in the formula
137
+
138
+
139
+
140
+ T
141
+ =
142
+
143
+
144
+ 1
145
+ 2
146
+
147
+
148
+ b
149
+ h
150
+
151
+
152
+ {\displaystyle T={\frac {1}{2}}bh}
153
+
154
+ derived above, the area of the triangle can be expressed as:
155
+
156
+ (where α is the interior angle at A, β is the interior angle at B,
157
+
158
+
159
+
160
+ γ
161
+
162
+
163
+ {\displaystyle \gamma }
164
+
165
+ is the interior angle at C and c is the line AB).
166
+
167
+ Furthermore, since sin α = sin (π − α) = sin (β +
168
+
169
+
170
+
171
+ γ
172
+
173
+
174
+ {\displaystyle \gamma }
175
+
176
+ ), and similarly for the other two angles:
177
+
178
+ Knowing AAS:
179
+
180
+ and analogously if the known side is a or c.
181
+
182
+ Knowing ASA:[12]
183
+
184
+ and analogously if the known side is b or c.
185
+
186
+ The shape of the triangle is determined by the lengths of the sides. Therefore, the area can also be derived from the lengths of the sides. By Heron's formula:
187
+
188
+ where
189
+
190
+
191
+
192
+ s
193
+ =
194
+
195
+
196
+
197
+
198
+ a
199
+ +
200
+ b
201
+ +
202
+ c
203
+
204
+ 2
205
+
206
+
207
+
208
+
209
+
210
+ {\displaystyle s={\tfrac {a+b+c}{2}}}
211
+
212
+ is the semiperimeter, or half of the triangle's perimeter.
213
+
214
+ Three other equivalent ways of writing Heron's formula are
215
+
216
+ The area of a parallelogram embedded in a three-dimensional Euclidean space can be calculated using vectors. Let vectors AB and AC point respectively from A to B and from A to C. The area of parallelogram ABDC is then
217
+
218
+ which is the magnitude of the cross product of vectors AB and AC. The area of triangle ABC is half of this,
219
+
220
+ The area of triangle ABC can also be expressed in terms of dot products as follows:
221
+
222
+ In two-dimensional Euclidean space, expressing vector AB as a free vector in Cartesian space equal to (x1,y1) and AC as (x2,y2), this can be rewritten as:
223
+
224
+ If vertex A is located at the origin (0, 0) of a Cartesian coordinate system and the coordinates of the other two vertices are given by B = (xB, yB) and C = (xC, yC), then the area can be computed as ​1⁄2 times the absolute value of the determinant
225
+
226
+ For three general vertices, the equation is:
227
+
228
+ which can be written as
229
+
230
+ If the points are labeled sequentially in the counterclockwise direction, the above determinant expressions are positive and the absolute value signs can be omitted.[13] The above formula is known as the shoelace formula or the surveyor's formula.
231
+
232
+ If we locate the vertices in the complex plane and denote them in counterclockwise sequence as a = xA + yAi, b = xB + yBi, and c = xC + yCi, and denote their complex conjugates as
233
+
234
+
235
+
236
+
237
+
238
+
239
+ a
240
+ ¯
241
+
242
+
243
+
244
+
245
+
246
+ {\displaystyle {\bar {a}}}
247
+
248
+ ,
249
+
250
+
251
+
252
+
253
+
254
+
255
+ b
256
+ ¯
257
+
258
+
259
+
260
+
261
+
262
+ {\displaystyle {\bar {b}}}
263
+
264
+ , and
265
+
266
+
267
+
268
+
269
+
270
+
271
+ c
272
+ ¯
273
+
274
+
275
+
276
+
277
+
278
+ {\displaystyle {\bar {c}}}
279
+
280
+ , then the formula
281
+
282
+ is equivalent to the shoelace formula.
283
+
284
+ In three dimensions, the area of a general triangle A = (xA, yA, zA), B = (xB, yB, zB) and C = (xC, yC, zC) is the Pythagorean sum of the areas of the respective projections on the three principal planes (i.e. x = 0, y = 0 and z = 0):
285
+
286
+ The area within any closed curve, such as a triangle, is given by the line integral around the curve of the algebraic or signed distance of a point on the curve from an arbitrary oriented straight line L. Points to the right of L as oriented are taken to be at negative distance from L, while the weight for the integral is taken to be the component of arc length parallel to L rather than arc length itself.
287
+
288
+ This method is well suited to computation of the area of an arbitrary polygon. Taking L to be the x-axis, the line integral between consecutive vertices (xi,yi) and (xi+1,yi+1) is given by the base times the mean height, namely (xi+1 − xi)(yi + yi+1)/2. The sign of the area is an overall indicator of the direction of traversal, with negative area indicating counterclockwise traversal. The area of a triangle then falls out as the case of a polygon with three sides.
289
+
290
+ While the line integral method has in common with other coordinate-based methods the arbitrary choice of a coordinate system, unlike the others it makes no arbitrary choice of vertex of the triangle as origin or of side as base. Furthermore, the choice of coordinate system defined by L commits to only two degrees of freedom rather than the usual three, since the weight is a local distance (e.g. xi+1 − xi in the above) whence the method does not require choosing an axis normal to L.
291
+
292
+ When working in polar coordinates it is not necessary to convert to Cartesian coordinates to use line integration, since the line integral between consecutive vertices (ri,θi) and (ri+1,θi+1) of a polygon is given directly by riri+1sin(θi+1 − θi)/2. This is valid for all values of θ, with some decrease in numerical accuracy when |θ| is many orders of magnitude greater than π. With this formulation negative area indicates clockwise traversal, which should be kept in mind when mixing polar and cartesian coordinates. Just as the choice of y-axis (x = 0) is immaterial for line integration in cartesian coordinates, so is the choice of zero heading (θ = 0) immaterial here.
293
+
294
+ Three formulas have the same structure as Heron's formula but are expressed in terms of different variables. First, denoting the medians from sides a, b, and c respectively as ma, mb, and mc and their semi-sum (ma + mb + mc)/2 as σ, we have[14]
295
+
296
+ Next, denoting the altitudes from sides a, b, and c respectively as ha, hb, and hc, and denoting the semi-sum of the reciprocals of the altitudes as
297
+
298
+
299
+
300
+ H
301
+ =
302
+ (
303
+
304
+ h
305
+
306
+ a
307
+
308
+
309
+ −
310
+ 1
311
+
312
+
313
+ +
314
+
315
+ h
316
+
317
+ b
318
+
319
+
320
+ −
321
+ 1
322
+
323
+
324
+ +
325
+
326
+ h
327
+
328
+ c
329
+
330
+
331
+ −
332
+ 1
333
+
334
+
335
+ )
336
+
337
+ /
338
+
339
+ 2
340
+
341
+
342
+ {\displaystyle H=(h_{a}^{-1}+h_{b}^{-1}+h_{c}^{-1})/2}
343
+
344
+ we have[15]
345
+
346
+ And denoting the semi-sum of the angles' sines as S = [(sin α) + (sin β) + (sin γ)]/2, we have[16]
347
+
348
+ where D is the diameter of the circumcircle:
349
+
350
+
351
+
352
+ D
353
+ =
354
+
355
+
356
+
357
+ a
358
+
359
+ sin
360
+ ⁡
361
+ α
362
+
363
+
364
+
365
+
366
+ =
367
+
368
+
369
+
370
+ b
371
+
372
+ sin
373
+ ⁡
374
+ β
375
+
376
+
377
+
378
+
379
+ =
380
+
381
+
382
+
383
+ c
384
+
385
+ sin
386
+ ⁡
387
+ γ
388
+
389
+
390
+
391
+
392
+ .
393
+
394
+
395
+ {\displaystyle D={\tfrac {a}{\sin \alpha }}={\tfrac {b}{\sin \beta }}={\tfrac {c}{\sin \gamma }}.}
396
+
397
+ See Pick's theorem for a technique for finding the area of any arbitrary lattice polygon (one drawn on a grid with vertically and horizontally adjacent lattice points at equal distances, and with vertices on lattice points).
398
+
399
+ The theorem states:
400
+
401
+ where
402
+
403
+
404
+
405
+ I
406
+
407
+
408
+ {\displaystyle I}
409
+
410
+ is the number of internal lattice points and B is the number of lattice points lying on the border of the polygon.
411
+
412
+ Numerous other area formulas exist, such as
413
+
414
+ where r is the inradius, and s is the semiperimeter (in fact, this formula holds for all tangential polygons), and[17]:Lemma 2
415
+
416
+ where
417
+
418
+
419
+
420
+
421
+ r
422
+
423
+ a
424
+
425
+
426
+ ,
427
+
428
+
429
+ r
430
+
431
+ b
432
+
433
+
434
+ ,
435
+
436
+
437
+ r
438
+
439
+ c
440
+
441
+
442
+
443
+
444
+ {\displaystyle r_{a},\,r_{b},\,r_{c}}
445
+
446
+ are the radii of the excircles tangent to sides a, b, c respectively.
447
+
448
+ We also have
449
+
450
+ and[18]
451
+
452
+ for circumdiameter D; and[19]
453
+
454
+ for angle α ≠ 90°.
455
+
456
+ The area can also be expressed as[20]
457
+
458
+ In 1885, Baker[21] gave a collection of over a hundred distinct area formulas for the triangle. These include:
459
+
460
+ for circumradius (radius of the circumcircle) R, and
461
+
462
+ The area T of any triangle with perimeter p satisfies
463
+
464
+ with equality holding if and only if the triangle is equilateral.[22][23]:657
465
+
466
+ Other upper bounds on the area T are given by[24]:p.290
467
+
468
+ and
469
+
470
+ both again holding if and only if the triangle is equilateral.
471
+
472
+ There are infinitely many lines that bisect the area of a triangle.[25] Three of them are the medians, which are the only area bisectors that go through the centroid. Three other area bisectors are parallel to the triangle's sides.
473
+
474
+ Any line through a triangle that splits both the triangle's area and its perimeter in half goes through the triangle's incenter. There can be one, two, or three of these for any given triangle.
475
+
476
+ The formulas in this section are true for all Euclidean triangles.
477
+
478
+ The medians and the sides are related by[26]:p.70
479
+
480
+ and
481
+
482
+ and equivalently for mb and mc.
483
+
484
+ For angle A opposite side a, the length of the internal angle bisector is given by[27]
485
+
486
+ for semiperimeter s, where the bisector length is measured from the vertex to where it meets the opposite side.
487
+
488
+ The interior perpendicular bisectors are given by
489
+
490
+ where the sides are
491
+
492
+
493
+
494
+ a
495
+ ≥
496
+ b
497
+ ≥
498
+ c
499
+
500
+
501
+ {\displaystyle a\geq b\geq c}
502
+
503
+ and the area is
504
+
505
+
506
+
507
+ T
508
+ .
509
+
510
+
511
+ {\displaystyle T.}
512
+
513
+ [28]:Thm 2
514
+
515
+ The altitude from, for example, the side of length a is
516
+
517
+ The following formulas involve the circumradius R and the inradius r:
518
+
519
+ where ha etc. are the altitudes to the subscripted sides;[26]:p.79
520
+
521
+ and
522
+
523
+ The product of two sides of a triangle equals the altitude to the third side times the diameter D of the circumcircle:[26]:p.64
524
+
525
+ Suppose two adjacent but non-overlapping triangles share the same side of length f and share the same circumcircle, so that the side of length f is a chord of the circumcircle and the triangles have side lengths (a, b, f) and (c, d, f), with the two triangles together forming a cyclic quadrilateral with side lengths in sequence (a, b, c, d). Then[29]:84
526
+
527
+ Let G be the centroid of a triangle with vertices A, B, and C, and let P be any interior point. Then the distances between the points are related by[29]:174
528
+
529
+ The sum of the squares of the triangle's sides equals three times the sum of the squared distances of the centroid from the vertices:
530
+
531
+ Let qa, qb, and qc be the distances from the centroid to the sides of lengths a, b, and c. Then[29]:173
532
+
533
+ and
534
+
535
+ for area T.
536
+
537
+ Carnot's theorem states that the sum of the distances from the circumcenter to the three sides equals the sum of the circumradius and the inradius.[26]:p.83 Here a segment's length is considered to be negative if and only if the segment lies entirely outside the triangle. This method is especially useful for deducing the properties of more abstract forms of triangles, such as the ones induced by Lie algebras, that otherwise have the same properties as usual triangles.
538
+
539
+ Euler's theorem states that the distance d between the circumcenter and the incenter is given by[26]:p.85
540
+
541
+ or equivalently
542
+
543
+ where R is the circumradius and r is the inradius. Thus for all triangles R ≥ 2r, with equality holding for equilateral triangles.
544
+
545
+ If we denote that the orthocenter divides one altitude into segments of lengths u and v, another altitude into segment lengths w and x, and the third altitude into segment lengths y and z, then uv = wx = yz.[26]:p.94
546
+
547
+ The distance from a side to the circumcenter equals half the distance from the opposite vertex to the orthocenter.[26]:p.99
548
+
549
+ The sum of the squares of the distances from the vertices to the orthocenter H plus the sum of the squares of the sides equals twelve times the square of the circumradius:[26]:p.102
550
+
551
+ In addition to the law of sines, the law of cosines, the law of tangents, and the trigonometric existence conditions given earlier, for any triangle
552
+
553
+ Morley's trisector theorem states that in any triangle, the three points of intersection of the adjacent angle trisectors form an equilateral triangle, called the Morley triangle.
554
+
555
+ As discussed above, every triangle has a unique inscribed circle (incircle) that is interior to the triangle and tangent to all three sides.
556
+
557
+ Every triangle has a unique Steiner inellipse which is interior to the triangle and tangent at the midpoints of the sides. Marden's theorem shows how to find the foci of this ellipse.[31] This ellipse has the greatest area of any ellipse tangent to all three sides of the triangle.
558
+
559
+ The Mandart inellipse of a triangle is the ellipse inscribed within the triangle tangent to its sides at the contact points of its excircles.
560
+
561
+ For any ellipse inscribed in a triangle ABC, let the foci be P and Q. Then[32]
562
+
563
+ Every convex polygon with area T can be inscribed in a triangle of area at most equal to 2T. Equality holds (exclusively) for a parallelogram.[33]
564
+
565
+ The Lemoine hexagon is a cyclic hexagon with vertices given by the six intersections of the sides of a triangle with the three lines that are parallel to the sides and that pass through its symmedian point. In either its simple form or its self-intersecting form, the Lemoine hexagon is interior to the triangle with two vertices on each side of the triangle.
566
+
567
+ Every acute triangle has three inscribed squares (squares in its interior such that all four of a square's vertices lie on a side of the triangle, so two of them lie on the same side and hence one side of the square coincides with part of a side of the triangle). In a right triangle two of the squares coincide and have a vertex at the triangle's right angle, so a right triangle has only two distinct inscribed squares. An obtuse triangle has only one inscribed square, with a side coinciding with part of the triangle's longest side. Within a given triangle, a longer common side is associated with a smaller inscribed square. If an inscribed square has side of length qa and the triangle has a side of length a, part of which side coincides with a side of the square, then qa, a, the altitude ha from the side a, and the triangle's area T are related according to[34][35]
568
+
569
+ The largest possible ratio of the area of the inscribed square to the area of the triangle is 1/2, which occurs when a2 = 2T, q = a/2, and the altitude of the triangle from the base of length a is equal to a. The smallest possible ratio of the side of one inscribed square to the side of another in the same non-obtuse triangle is
570
+
571
+
572
+
573
+ 2
574
+
575
+
576
+ 2
577
+
578
+
579
+
580
+ /
581
+
582
+ 3
583
+ =
584
+ 0.94....
585
+
586
+
587
+ {\displaystyle 2{\sqrt {2}}/3=0.94....}
588
+
589
+ [35] Both of these extreme cases occur for the isosceles right triangle.
590
+
591
+ From an interior point in a reference triangle, the nearest points on the three sides serve as the vertices of the pedal triangle of that point. If the interior point is the circumcenter of the reference triangle, the vertices of the pedal triangle are the midpoints of the reference triangle's sides, and so the pedal triangle is called the midpoint triangle or medial triangle. The midpoint triangle subdivides the reference triangle into four congruent triangles which are similar to the reference triangle.
592
+
593
+ The Gergonne triangle or intouch triangle of a reference triangle has its vertices at the three points of tangency of the reference triangle's sides with its incircle. The extouch triangle of a reference triangle has its vertices at the points of tangency of the reference triangle's excircles with its sides (not extended).
594
+
595
+ The tangential triangle of a reference triangle (other than a right triangle) is the triangle whose sides are on the tangent lines to the reference triangle's circumcircle at its vertices.
596
+
597
+ As mentioned above, every triangle has a unique circumcircle, a circle passing through all three vertices, whose center is the intersection of the perpendicular bisectors of the triangle's sides.
598
+
599
+ Further, every triangle has a unique Steiner circumellipse, which passes through the triangle's vertices and has its center at the triangle's centroid. Of all ellipses going through the triangle's vertices, it has the smallest area.
600
+
601
+ The Kiepert hyperbola is the unique conic which passes through the triangle's three vertices, its centroid, and its circumcenter.
602
+
603
+ Of all triangles contained in a given convex polygon, there exists a triangle with maximal area whose vertices are all vertices of the given polygon.[36]
604
+
605
+ One way to identify locations of points in (or outside) a triangle is to place the triangle in an arbitrary location and orientation in the Cartesian plane, and to use Cartesian coordinates. While convenient for many purposes, this approach has the disadvantage of all points' coordinate values being dependent on the arbitrary placement in the plane.
606
+
607
+ Two systems avoid that feature, so that the coordinates of a point are not affected by moving the triangle, rotating it, or reflecting it as in a mirror, any of which give a congruent triangle, or even by rescaling it to give a similar triangle:
608
+
609
+ A non-planar triangle is a triangle which is not contained in a (flat) plane. Some examples of non-planar triangles in non-Euclidean geometries are spherical triangles in spherical geometry and hyperbolic triangles in hyperbolic geometry.
610
+
611
+ While the measures of the internal angles in planar triangles always sum to 180°, a hyperbolic triangle has measures of angles that sum to less than 180°, and a spherical triangle has measures of angles that sum to more than 180°. A hyperbolic triangle can be obtained by drawing on a negatively curved surface, such as a saddle surface, and a spherical triangle can be obtained by drawing on a positively curved surface such as a sphere. Thus, if one draws a giant triangle on the surface of the Earth, one will find that the sum of the measures of its angles is greater than 180°; in fact it will be between 180° and 540°.[37] In particular it is possible to draw a triangle on a sphere such that the measure of each of its internal angles is equal to 90°, adding up to a total of 270°.
612
+
613
+ Specifically, on a sphere the sum of the angles of a triangle is
614
+
615
+ where f is the fraction of the sphere's area which is enclosed by the triangle. For example, suppose that we draw a triangle on the Earth's surface with vertices at the North Pole, at a point on the equator at 0° longitude, and a point on the equator at 90° West longitude. The great circle line between the latter two points is the equator, and the great circle line between either of those points and the North Pole is a line of longitude; so there are right angles at the two points on the equator. Moreover, the angle at the North Pole is also 90° because the other two vertices differ by 90° of longitude. So the sum of the angles in this triangle is 90° + 90° + 90° = 270°. The triangle encloses 1/4 of the northern hemisphere (90°/360° as viewed from the North Pole) and therefore 1/8 of the Earth's surface, so in the formula f = 1/8; thus the formula correctly gives the sum of the triangle's angles as 270°.
616
+
617
+ From the above angle sum formula we can also see that the Earth's surface is locally flat: If we draw an arbitrarily small triangle in the neighborhood of one point on the Earth's surface, the fraction f of the Earth's surface which is enclosed by the triangle will be arbitrarily close to zero. In this case the angle sum formula simplifies to 180°, which we know is what Euclidean geometry tells us for triangles on a flat surface.
618
+
619
+ Rectangles have been the most popular and common geometric form for buildings since the shape is easy to stack and organize; as a standard, it is easy to design furniture and fixtures to fit inside rectangularly shaped buildings. But triangles, while more difficult to use conceptually, provide a great deal of strength. As computer technology helps architects design creative new buildings, triangular shapes are becoming increasingly prevalent as parts of buildings and as the primary shape for some types of skyscrapers as well as building materials. In Tokyo in 1989, architects had wondered whether it was possible to build a 500-story tower to provide affordable office space for this densely packed city, but with the danger to buildings from earthquakes, architects considered that a triangular shape would be necessary if such a building were to be built.[38]
620
+
621
+ In New York City, as Broadway crisscrosses major avenues, the resulting blocks are cut like triangles, and buildings have been built on these shapes; one such building is the triangularly shaped Flatiron Building which real estate people admit has a "warren of awkward spaces that do not easily accommodate modern office furniture" but that has not prevented the structure from becoming a landmark icon.[39] Designers have made houses in Norway using triangular themes.[40] Triangle shapes have appeared in churches[41] as well as public buildings including colleges[42] as well as supports for innovative home designs.[43]
622
+
623
+ Triangles are sturdy; while a rectangle can collapse into a parallelogram from pressure to one of its points, triangles have a natural strength which supports structures against lateral pressures. A triangle will not change shape unless its sides are bent or extended or broken or if its joints break; in essence, each of the three sides supports the other two. A rectangle, in contrast, is more dependent on the strength of its joints in a structural sense. Some innovative designers have proposed making bricks not out of rectangles, but with triangular shapes which can be combined in three dimensions.[44] It is likely that triangles will be used increasingly in new ways as architecture increases in complexity. It is important to remember that triangles are strong in terms of rigidity, but while packed in a tessellating arrangement triangles are not as strong as hexagons under compression (hence the prevalence of hexagonal forms in nature). Tessellated triangles still maintain superior strength for cantilevering however, and this is the basis for one of the strongest man made structures, the tetrahedral truss.
624
+
625
+
626
+
627
+
628
+
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1
+
2
+
3
+ A triangle is a polygon with three edges and three vertices. It is one of the basic shapes in geometry. A triangle with vertices A, B, and C is denoted
4
+
5
+
6
+
7
+ △
8
+ A
9
+ B
10
+ C
11
+
12
+
13
+ {\displaystyle \triangle ABC}
14
+
15
+ .
16
+
17
+ In Euclidean geometry any three points, when non-collinear, determine a unique triangle and simultaneously, a unique plane (i.e. a two-dimensional Euclidean space). In other words, there is only one plane that contains that triangle, and every triangle is contained in some plane. If the entire geometry is only the Euclidean plane, there is only one plane and all triangles are contained in it; however, in higher-dimensional Euclidean spaces, this is no longer true. This article is about triangles in Euclidean geometry, and in particular, the Euclidean plane, except where otherwise noted.
18
+
19
+ Triangles can be classified according to the lengths of their sides:
20
+
21
+ Hatch marks, also called tick marks, are used in diagrams of triangles and other geometric figures to identify sides of equal lengths. A side can be marked with a pattern of "ticks", short line segments in the form of tally marks; two sides have equal lengths if they are both marked with the same pattern. In a triangle, the pattern is usually no more than 3 ticks. An equilateral triangle has the same pattern on all 3 sides, an isosceles triangle has the same pattern on just 2 sides, and a scalene triangle has different patterns on all sides since no sides are equal. Similarly, patterns of 1, 2, or 3 concentric arcs inside the angles are used to indicate equal angles. An equilateral triangle has the same pattern on all 3 angles, an isosceles triangle has the same pattern on just 2 angles, and a scalene triangle has different patterns on all angles since no angles are equal.
22
+
23
+ Triangles can also be classified according to their internal angles, measured here in degrees.
24
+
25
+ A triangle that has two angles with the same measure also has two sides with the same length, and therefore it is an isosceles triangle. It follows that in a triangle where all angles have the same measure, all three sides have the same length, and such a triangle is therefore equilateral.
26
+
27
+ Triangles are assumed to be two-dimensional plane figures, unless the context provides otherwise (see Non-planar triangles, below). In rigorous treatments, a triangle is therefore called a 2-simplex (see also Polytope). Elementary facts about triangles were presented by Euclid in books 1–4 of his Elements, around 300 BC.
28
+
29
+ The sum of the measures of the interior angles of a triangle in Euclidean space is always 180 degrees.[5] This fact is equivalent to Euclid's parallel postulate. This allows determination of the measure of the third angle of any triangle given the measure of two angles. An exterior angle of a triangle is an angle that is a linear pair (and hence supplementary) to an interior angle. The measure of an exterior angle of a triangle is equal to the sum of the measures of the two interior angles that are not adjacent to it; this is the exterior angle theorem. The sum of the measures of the three exterior angles (one for each vertex) of any triangle is 360 degrees.[note 2]
30
+
31
+ Two triangles are said to be similar if every angle of one triangle has the same measure as the corresponding angle in the other triangle. The corresponding sides of similar triangles have lengths that are in the same proportion, and this property is also sufficient to establish similarity.
32
+
33
+ Some basic theorems about similar triangles are:
34
+
35
+ Two triangles that are congruent have exactly the same size and shape:[note 4] all pairs of corresponding interior angles are equal in measure, and all pairs of corresponding sides have the same length. (This is a total of six equalities, but three are often sufficient to prove congruence.)
36
+
37
+ Some individually necessary and sufficient conditions for a pair of triangles to be congruent are:
38
+
39
+ Some individually sufficient conditions are:
40
+
41
+ An important condition is:
42
+
43
+ Using right triangles and the concept of similarity, the trigonometric functions sine and cosine can be defined. These are functions of an angle which are investigated in trigonometry.
44
+
45
+ A central theorem is the Pythagorean theorem, which states in any right triangle, the square of the length of the hypotenuse equals the sum of the squares of the lengths of the two other sides. If the hypotenuse has length c, and the legs have lengths a and b, then the theorem states that
46
+
47
+ The converse is true: if the lengths of the sides of a triangle satisfy the above equation, then the triangle has a right angle opposite side c.
48
+
49
+ Some other facts about right triangles:
50
+
51
+ For all triangles, angles and sides are related by the law of cosines and law of sines (also called the cosine rule and sine rule).
52
+
53
+ The triangle inequality states that the sum of the lengths of any two sides of a triangle must be greater than or equal to the length of the third side. That sum can equal the length of the third side only in the case of a degenerate triangle, one with collinear vertices. It is not possible for that sum to be less than the length of the third side. A triangle with three given positive side lengths exists if and only if those side lengths satisfy the triangle inequality.
54
+
55
+ Three given angles form a non-degenerate triangle (and indeed an infinitude of them) if and only if both of these conditions hold: (a) each of the angles is positive, and (b) the angles sum to 180°. If degenerate triangles are permitted, angles of 0° are permitted.
56
+
57
+ Three positive angles α, β, and γ, each of them less than 180°, are the angles of a triangle if and only if any one of the following conditions holds:
58
+
59
+ the last equality applying only if none of the angles is 90° (so the tangent function's value is always finite).
60
+
61
+ There are thousands of different constructions that find a special point associated with (and often inside) a triangle, satisfying some unique property: see the article Encyclopedia of Triangle Centers for a catalogue of them. Often they are constructed by finding three lines associated in a symmetrical way with the three sides (or vertices) and then proving that the three lines meet in a single point: an important tool for proving the existence of these is Ceva's theorem, which gives a criterion for determining when three such lines are concurrent. Similarly, lines associated with a triangle are often constructed by proving that three symmetrically constructed points are collinear: here Menelaus' theorem gives a useful general criterion. In this section just a few of the most commonly encountered constructions are explained.
62
+
63
+ A perpendicular bisector of a side of a triangle is a straight line passing through the midpoint of the side and being perpendicular to it, i.e. forming a right angle with it. The three perpendicular bisectors meet in a single point, the triangle's circumcenter, usually denoted by O; this point is the center of the circumcircle, the circle passing through all three vertices. The diameter of this circle, called the circumdiameter, can be found from the law of sines stated above. The circumcircle's radius is called the circumradius.
64
+
65
+ Thales' theorem implies that if the circumcenter is located on a side of the triangle, then the opposite angle is a right one. If the circumcenter is located inside the triangle, then the triangle is acute; if the circumcenter is located outside the triangle, then the triangle is obtuse.
66
+
67
+ An altitude of a triangle is a straight line through a vertex and perpendicular to (i.e. forming a right angle with) the opposite side. This opposite side is called the base of the altitude, and the point where the altitude intersects the base (or its extension) is called the foot of the altitude. The length of the altitude is the distance between the base and the vertex. The three altitudes intersect in a single point, called the orthocenter of the triangle, usually denoted by H. The orthocenter lies inside the triangle if and only if the triangle is acute.
68
+
69
+ An angle bisector of a triangle is a straight line through a vertex which cuts the corresponding angle in half. The three angle bisectors intersect in a single point, the incenter, usually denoted by I, the center of the triangle's incircle. The incircle is the circle which lies inside the triangle and touches all three sides. Its radius is called the inradius. There are three other important circles, the excircles; they lie outside the triangle and touch one side as well as the extensions of the other two. The centers of the in- and excircles form an orthocentric system.
70
+
71
+ A median of a triangle is a straight line through a vertex and the midpoint of the opposite side, and divides the triangle into two equal areas. The three medians intersect in a single point, the triangle's centroid or geometric barycenter, usually denoted by G. The centroid of a rigid triangular object (cut out of a thin sheet of uniform density) is also its center of mass: the object can be balanced on its centroid in a uniform gravitational field. The centroid cuts every median in the ratio 2:1, i.e. the distance between a vertex and the centroid is twice the distance between the centroid and the midpoint of the opposite side.
72
+
73
+ The midpoints of the three sides and the feet of the three altitudes all lie on a single circle, the triangle's nine-point circle. The remaining three points for which it is named are the midpoints of the portion of altitude between the vertices and the orthocenter. The radius of the nine-point circle is half that of the circumcircle. It touches the incircle (at the Feuerbach point) and the three excircles.
74
+
75
+ The orthocenter (blue point), center of the nine-point circle (red), centroid (orange), and circumcenter (green) all lie on a single line, known as Euler's line (red line). The center of the nine-point circle lies at the midpoint between the orthocenter and the circumcenter, and the distance between the centroid and the circumcenter is half that between the centroid and the orthocenter.
76
+
77
+ The center of the incircle is not in general located on Euler's line.
78
+
79
+ If one reflects a median in the angle bisector that passes through the same vertex, one obtains a symmedian. The three symmedians intersect in a single point, the symmedian point of the triangle.
80
+
81
+ There are various standard methods for calculating the length of a side or the measure of an angle. Certain methods are suited to calculating values in a right-angled triangle; more complex methods may be required in other situations.
82
+
83
+ In right triangles, the trigonometric ratios of sine, cosine and tangent can be used to find unknown angles and the lengths of unknown sides. The sides of the triangle are known as follows:
84
+
85
+ The sine of an angle is the ratio of the length of the opposite side to the length of the hypotenuse. In our case
86
+
87
+ This ratio does not depend on the particular right triangle chosen, as long as it contains the angle A, since all those triangles are similar.
88
+
89
+ The cosine of an angle is the ratio of the length of the adjacent side to the length of the hypotenuse. In our case
90
+
91
+ The tangent of an angle is the ratio of the length of the opposite side to the length of the adjacent side. In our case
92
+
93
+ The acronym "SOH-CAH-TOA" is a useful mnemonic for these ratios.
94
+
95
+ The inverse trigonometric functions can be used to calculate the internal angles for a right angled triangle with the length of any two sides.
96
+
97
+ Arcsin can be used to calculate an angle from the length of the opposite side and the length of the hypotenuse.
98
+
99
+ Arccos can be used to calculate an angle from the length of the adjacent side and the length of the hypotenuse.
100
+
101
+ Arctan can be used to calculate an angle from the length of the opposite side and the length of the adjacent side.
102
+
103
+ In introductory geometry and trigonometry courses, the notation sin−1, cos−1, etc., are often used in place of arcsin, arccos, etc. However, the arcsin, arccos, etc., notation is standard in higher mathematics where trigonometric functions are commonly raised to powers, as this avoids confusion between multiplicative inverse and compositional inverse.
104
+
105
+ The law of sines, or sine rule,[8] states that the ratio of the length of a side to the sine of its corresponding opposite angle is constant, that is
106
+
107
+ This ratio is equal to the diameter of the circumscribed circle of the given triangle. Another interpretation of this theorem is that every triangle with angles α, β and γ is similar to a triangle with side lengths equal to sin α, sin β and sin γ. This triangle can be constructed by first constructing a circle of diameter 1, and inscribing in it two of the angles of the triangle. The length of the sides of that triangle will be sin α, sin β and sin γ. The side whose length is sin α is opposite to the angle whose measure is α, etc.
108
+
109
+ The law of cosines, or cosine rule, connects the length of an unknown side of a triangle to the length of the other sides and the angle opposite to the unknown side.[8] As per the law:
110
+
111
+ For a triangle with length of sides a, b, c and angles of α, β, γ respectively, given two known lengths of a triangle a and b, and the angle between the two known sides γ (or the angle opposite to the unknown side c), to calculate the third side c, the following formula can be used:
112
+
113
+ If the lengths of all three sides of any triangle are known the three angles can be calculated:
114
+
115
+ The law of tangents, or tangent rule, can be used to find a side or an angle when two sides and an angle or two angles and a side are known. It states that:[9]
116
+
117
+ "Solution of triangles" is the main trigonometric problem: to find missing characteristics of a triangle (three angles, the lengths of the three sides etc.) when at least three of these characteristics are given. The triangle can be located on a plane or on a sphere. This problem often occurs in various trigonometric applications, such as geodesy, astronomy, construction, navigation etc.
118
+
119
+ Calculating the area T of a triangle is an elementary problem encountered often in many different situations. The best known and simplest formula is:
120
+
121
+ where b is the length of the base of the triangle, and h is the height or altitude of the triangle. The term "base" denotes any side, and "height" denotes the length of a perpendicular from the vertex opposite the base onto the line containing the base. In 499 CE Aryabhata, used this illustrated method in the Aryabhatiya (section 2.6).[10]
122
+
123
+ Although simple, this formula is only useful if the height can be readily found, which is not always the case. For example, the surveyor of a triangular field might find it relatively easy to measure the length of each side, but relatively difficult to construct a 'height'. Various methods may be used in practice, depending on what is known about the triangle. The following is a selection of frequently used formulae for the area of a triangle.[11]
124
+
125
+ The height of a triangle can be found through the application of trigonometry.
126
+
127
+ Knowing SAS: Using the labels in the image on the right, the altitude is h = a sin
128
+
129
+
130
+
131
+ γ
132
+
133
+
134
+ {\displaystyle \gamma }
135
+
136
+ . Substituting this in the formula
137
+
138
+
139
+
140
+ T
141
+ =
142
+
143
+
144
+ 1
145
+ 2
146
+
147
+
148
+ b
149
+ h
150
+
151
+
152
+ {\displaystyle T={\frac {1}{2}}bh}
153
+
154
+ derived above, the area of the triangle can be expressed as:
155
+
156
+ (where α is the interior angle at A, β is the interior angle at B,
157
+
158
+
159
+
160
+ γ
161
+
162
+
163
+ {\displaystyle \gamma }
164
+
165
+ is the interior angle at C and c is the line AB).
166
+
167
+ Furthermore, since sin α = sin (π − α) = sin (β +
168
+
169
+
170
+
171
+ γ
172
+
173
+
174
+ {\displaystyle \gamma }
175
+
176
+ ), and similarly for the other two angles:
177
+
178
+ Knowing AAS:
179
+
180
+ and analogously if the known side is a or c.
181
+
182
+ Knowing ASA:[12]
183
+
184
+ and analogously if the known side is b or c.
185
+
186
+ The shape of the triangle is determined by the lengths of the sides. Therefore, the area can also be derived from the lengths of the sides. By Heron's formula:
187
+
188
+ where
189
+
190
+
191
+
192
+ s
193
+ =
194
+
195
+
196
+
197
+
198
+ a
199
+ +
200
+ b
201
+ +
202
+ c
203
+
204
+ 2
205
+
206
+
207
+
208
+
209
+
210
+ {\displaystyle s={\tfrac {a+b+c}{2}}}
211
+
212
+ is the semiperimeter, or half of the triangle's perimeter.
213
+
214
+ Three other equivalent ways of writing Heron's formula are
215
+
216
+ The area of a parallelogram embedded in a three-dimensional Euclidean space can be calculated using vectors. Let vectors AB and AC point respectively from A to B and from A to C. The area of parallelogram ABDC is then
217
+
218
+ which is the magnitude of the cross product of vectors AB and AC. The area of triangle ABC is half of this,
219
+
220
+ The area of triangle ABC can also be expressed in terms of dot products as follows:
221
+
222
+ In two-dimensional Euclidean space, expressing vector AB as a free vector in Cartesian space equal to (x1,y1) and AC as (x2,y2), this can be rewritten as:
223
+
224
+ If vertex A is located at the origin (0, 0) of a Cartesian coordinate system and the coordinates of the other two vertices are given by B = (xB, yB) and C = (xC, yC), then the area can be computed as ​1⁄2 times the absolute value of the determinant
225
+
226
+ For three general vertices, the equation is:
227
+
228
+ which can be written as
229
+
230
+ If the points are labeled sequentially in the counterclockwise direction, the above determinant expressions are positive and the absolute value signs can be omitted.[13] The above formula is known as the shoelace formula or the surveyor's formula.
231
+
232
+ If we locate the vertices in the complex plane and denote them in counterclockwise sequence as a = xA + yAi, b = xB + yBi, and c = xC + yCi, and denote their complex conjugates as
233
+
234
+
235
+
236
+
237
+
238
+
239
+ a
240
+ ¯
241
+
242
+
243
+
244
+
245
+
246
+ {\displaystyle {\bar {a}}}
247
+
248
+ ,
249
+
250
+
251
+
252
+
253
+
254
+
255
+ b
256
+ ¯
257
+
258
+
259
+
260
+
261
+
262
+ {\displaystyle {\bar {b}}}
263
+
264
+ , and
265
+
266
+
267
+
268
+
269
+
270
+
271
+ c
272
+ ¯
273
+
274
+
275
+
276
+
277
+
278
+ {\displaystyle {\bar {c}}}
279
+
280
+ , then the formula
281
+
282
+ is equivalent to the shoelace formula.
283
+
284
+ In three dimensions, the area of a general triangle A = (xA, yA, zA), B = (xB, yB, zB) and C = (xC, yC, zC) is the Pythagorean sum of the areas of the respective projections on the three principal planes (i.e. x = 0, y = 0 and z = 0):
285
+
286
+ The area within any closed curve, such as a triangle, is given by the line integral around the curve of the algebraic or signed distance of a point on the curve from an arbitrary oriented straight line L. Points to the right of L as oriented are taken to be at negative distance from L, while the weight for the integral is taken to be the component of arc length parallel to L rather than arc length itself.
287
+
288
+ This method is well suited to computation of the area of an arbitrary polygon. Taking L to be the x-axis, the line integral between consecutive vertices (xi,yi) and (xi+1,yi+1) is given by the base times the mean height, namely (xi+1 − xi)(yi + yi+1)/2. The sign of the area is an overall indicator of the direction of traversal, with negative area indicating counterclockwise traversal. The area of a triangle then falls out as the case of a polygon with three sides.
289
+
290
+ While the line integral method has in common with other coordinate-based methods the arbitrary choice of a coordinate system, unlike the others it makes no arbitrary choice of vertex of the triangle as origin or of side as base. Furthermore, the choice of coordinate system defined by L commits to only two degrees of freedom rather than the usual three, since the weight is a local distance (e.g. xi+1 − xi in the above) whence the method does not require choosing an axis normal to L.
291
+
292
+ When working in polar coordinates it is not necessary to convert to Cartesian coordinates to use line integration, since the line integral between consecutive vertices (ri,θi) and (ri+1,θi+1) of a polygon is given directly by riri+1sin(θi+1 − θi)/2. This is valid for all values of θ, with some decrease in numerical accuracy when |θ| is many orders of magnitude greater than π. With this formulation negative area indicates clockwise traversal, which should be kept in mind when mixing polar and cartesian coordinates. Just as the choice of y-axis (x = 0) is immaterial for line integration in cartesian coordinates, so is the choice of zero heading (θ = 0) immaterial here.
293
+
294
+ Three formulas have the same structure as Heron's formula but are expressed in terms of different variables. First, denoting the medians from sides a, b, and c respectively as ma, mb, and mc and their semi-sum (ma + mb + mc)/2 as σ, we have[14]
295
+
296
+ Next, denoting the altitudes from sides a, b, and c respectively as ha, hb, and hc, and denoting the semi-sum of the reciprocals of the altitudes as
297
+
298
+
299
+
300
+ H
301
+ =
302
+ (
303
+
304
+ h
305
+
306
+ a
307
+
308
+
309
+ −
310
+ 1
311
+
312
+
313
+ +
314
+
315
+ h
316
+
317
+ b
318
+
319
+
320
+ −
321
+ 1
322
+
323
+
324
+ +
325
+
326
+ h
327
+
328
+ c
329
+
330
+
331
+ −
332
+ 1
333
+
334
+
335
+ )
336
+
337
+ /
338
+
339
+ 2
340
+
341
+
342
+ {\displaystyle H=(h_{a}^{-1}+h_{b}^{-1}+h_{c}^{-1})/2}
343
+
344
+ we have[15]
345
+
346
+ And denoting the semi-sum of the angles' sines as S = [(sin α) + (sin β) + (sin γ)]/2, we have[16]
347
+
348
+ where D is the diameter of the circumcircle:
349
+
350
+
351
+
352
+ D
353
+ =
354
+
355
+
356
+
357
+ a
358
+
359
+ sin
360
+ ⁡
361
+ α
362
+
363
+
364
+
365
+
366
+ =
367
+
368
+
369
+
370
+ b
371
+
372
+ sin
373
+ ⁡
374
+ β
375
+
376
+
377
+
378
+
379
+ =
380
+
381
+
382
+
383
+ c
384
+
385
+ sin
386
+ ⁡
387
+ γ
388
+
389
+
390
+
391
+
392
+ .
393
+
394
+
395
+ {\displaystyle D={\tfrac {a}{\sin \alpha }}={\tfrac {b}{\sin \beta }}={\tfrac {c}{\sin \gamma }}.}
396
+
397
+ See Pick's theorem for a technique for finding the area of any arbitrary lattice polygon (one drawn on a grid with vertically and horizontally adjacent lattice points at equal distances, and with vertices on lattice points).
398
+
399
+ The theorem states:
400
+
401
+ where
402
+
403
+
404
+
405
+ I
406
+
407
+
408
+ {\displaystyle I}
409
+
410
+ is the number of internal lattice points and B is the number of lattice points lying on the border of the polygon.
411
+
412
+ Numerous other area formulas exist, such as
413
+
414
+ where r is the inradius, and s is the semiperimeter (in fact, this formula holds for all tangential polygons), and[17]:Lemma 2
415
+
416
+ where
417
+
418
+
419
+
420
+
421
+ r
422
+
423
+ a
424
+
425
+
426
+ ,
427
+
428
+
429
+ r
430
+
431
+ b
432
+
433
+
434
+ ,
435
+
436
+
437
+ r
438
+
439
+ c
440
+
441
+
442
+
443
+
444
+ {\displaystyle r_{a},\,r_{b},\,r_{c}}
445
+
446
+ are the radii of the excircles tangent to sides a, b, c respectively.
447
+
448
+ We also have
449
+
450
+ and[18]
451
+
452
+ for circumdiameter D; and[19]
453
+
454
+ for angle α ≠ 90°.
455
+
456
+ The area can also be expressed as[20]
457
+
458
+ In 1885, Baker[21] gave a collection of over a hundred distinct area formulas for the triangle. These include:
459
+
460
+ for circumradius (radius of the circumcircle) R, and
461
+
462
+ The area T of any triangle with perimeter p satisfies
463
+
464
+ with equality holding if and only if the triangle is equilateral.[22][23]:657
465
+
466
+ Other upper bounds on the area T are given by[24]:p.290
467
+
468
+ and
469
+
470
+ both again holding if and only if the triangle is equilateral.
471
+
472
+ There are infinitely many lines that bisect the area of a triangle.[25] Three of them are the medians, which are the only area bisectors that go through the centroid. Three other area bisectors are parallel to the triangle's sides.
473
+
474
+ Any line through a triangle that splits both the triangle's area and its perimeter in half goes through the triangle's incenter. There can be one, two, or three of these for any given triangle.
475
+
476
+ The formulas in this section are true for all Euclidean triangles.
477
+
478
+ The medians and the sides are related by[26]:p.70
479
+
480
+ and
481
+
482
+ and equivalently for mb and mc.
483
+
484
+ For angle A opposite side a, the length of the internal angle bisector is given by[27]
485
+
486
+ for semiperimeter s, where the bisector length is measured from the vertex to where it meets the opposite side.
487
+
488
+ The interior perpendicular bisectors are given by
489
+
490
+ where the sides are
491
+
492
+
493
+
494
+ a
495
+ ≥
496
+ b
497
+ ≥
498
+ c
499
+
500
+
501
+ {\displaystyle a\geq b\geq c}
502
+
503
+ and the area is
504
+
505
+
506
+
507
+ T
508
+ .
509
+
510
+
511
+ {\displaystyle T.}
512
+
513
+ [28]:Thm 2
514
+
515
+ The altitude from, for example, the side of length a is
516
+
517
+ The following formulas involve the circumradius R and the inradius r:
518
+
519
+ where ha etc. are the altitudes to the subscripted sides;[26]:p.79
520
+
521
+ and
522
+
523
+ The product of two sides of a triangle equals the altitude to the third side times the diameter D of the circumcircle:[26]:p.64
524
+
525
+ Suppose two adjacent but non-overlapping triangles share the same side of length f and share the same circumcircle, so that the side of length f is a chord of the circumcircle and the triangles have side lengths (a, b, f) and (c, d, f), with the two triangles together forming a cyclic quadrilateral with side lengths in sequence (a, b, c, d). Then[29]:84
526
+
527
+ Let G be the centroid of a triangle with vertices A, B, and C, and let P be any interior point. Then the distances between the points are related by[29]:174
528
+
529
+ The sum of the squares of the triangle's sides equals three times the sum of the squared distances of the centroid from the vertices:
530
+
531
+ Let qa, qb, and qc be the distances from the centroid to the sides of lengths a, b, and c. Then[29]:173
532
+
533
+ and
534
+
535
+ for area T.
536
+
537
+ Carnot's theorem states that the sum of the distances from the circumcenter to the three sides equals the sum of the circumradius and the inradius.[26]:p.83 Here a segment's length is considered to be negative if and only if the segment lies entirely outside the triangle. This method is especially useful for deducing the properties of more abstract forms of triangles, such as the ones induced by Lie algebras, that otherwise have the same properties as usual triangles.
538
+
539
+ Euler's theorem states that the distance d between the circumcenter and the incenter is given by[26]:p.85
540
+
541
+ or equivalently
542
+
543
+ where R is the circumradius and r is the inradius. Thus for all triangles R ≥ 2r, with equality holding for equilateral triangles.
544
+
545
+ If we denote that the orthocenter divides one altitude into segments of lengths u and v, another altitude into segment lengths w and x, and the third altitude into segment lengths y and z, then uv = wx = yz.[26]:p.94
546
+
547
+ The distance from a side to the circumcenter equals half the distance from the opposite vertex to the orthocenter.[26]:p.99
548
+
549
+ The sum of the squares of the distances from the vertices to the orthocenter H plus the sum of the squares of the sides equals twelve times the square of the circumradius:[26]:p.102
550
+
551
+ In addition to the law of sines, the law of cosines, the law of tangents, and the trigonometric existence conditions given earlier, for any triangle
552
+
553
+ Morley's trisector theorem states that in any triangle, the three points of intersection of the adjacent angle trisectors form an equilateral triangle, called the Morley triangle.
554
+
555
+ As discussed above, every triangle has a unique inscribed circle (incircle) that is interior to the triangle and tangent to all three sides.
556
+
557
+ Every triangle has a unique Steiner inellipse which is interior to the triangle and tangent at the midpoints of the sides. Marden's theorem shows how to find the foci of this ellipse.[31] This ellipse has the greatest area of any ellipse tangent to all three sides of the triangle.
558
+
559
+ The Mandart inellipse of a triangle is the ellipse inscribed within the triangle tangent to its sides at the contact points of its excircles.
560
+
561
+ For any ellipse inscribed in a triangle ABC, let the foci be P and Q. Then[32]
562
+
563
+ Every convex polygon with area T can be inscribed in a triangle of area at most equal to 2T. Equality holds (exclusively) for a parallelogram.[33]
564
+
565
+ The Lemoine hexagon is a cyclic hexagon with vertices given by the six intersections of the sides of a triangle with the three lines that are parallel to the sides and that pass through its symmedian point. In either its simple form or its self-intersecting form, the Lemoine hexagon is interior to the triangle with two vertices on each side of the triangle.
566
+
567
+ Every acute triangle has three inscribed squares (squares in its interior such that all four of a square's vertices lie on a side of the triangle, so two of them lie on the same side and hence one side of the square coincides with part of a side of the triangle). In a right triangle two of the squares coincide and have a vertex at the triangle's right angle, so a right triangle has only two distinct inscribed squares. An obtuse triangle has only one inscribed square, with a side coinciding with part of the triangle's longest side. Within a given triangle, a longer common side is associated with a smaller inscribed square. If an inscribed square has side of length qa and the triangle has a side of length a, part of which side coincides with a side of the square, then qa, a, the altitude ha from the side a, and the triangle's area T are related according to[34][35]
568
+
569
+ The largest possible ratio of the area of the inscribed square to the area of the triangle is 1/2, which occurs when a2 = 2T, q = a/2, and the altitude of the triangle from the base of length a is equal to a. The smallest possible ratio of the side of one inscribed square to the side of another in the same non-obtuse triangle is
570
+
571
+
572
+
573
+ 2
574
+
575
+
576
+ 2
577
+
578
+
579
+
580
+ /
581
+
582
+ 3
583
+ =
584
+ 0.94....
585
+
586
+
587
+ {\displaystyle 2{\sqrt {2}}/3=0.94....}
588
+
589
+ [35] Both of these extreme cases occur for the isosceles right triangle.
590
+
591
+ From an interior point in a reference triangle, the nearest points on the three sides serve as the vertices of the pedal triangle of that point. If the interior point is the circumcenter of the reference triangle, the vertices of the pedal triangle are the midpoints of the reference triangle's sides, and so the pedal triangle is called the midpoint triangle or medial triangle. The midpoint triangle subdivides the reference triangle into four congruent triangles which are similar to the reference triangle.
592
+
593
+ The Gergonne triangle or intouch triangle of a reference triangle has its vertices at the three points of tangency of the reference triangle's sides with its incircle. The extouch triangle of a reference triangle has its vertices at the points of tangency of the reference triangle's excircles with its sides (not extended).
594
+
595
+ The tangential triangle of a reference triangle (other than a right triangle) is the triangle whose sides are on the tangent lines to the reference triangle's circumcircle at its vertices.
596
+
597
+ As mentioned above, every triangle has a unique circumcircle, a circle passing through all three vertices, whose center is the intersection of the perpendicular bisectors of the triangle's sides.
598
+
599
+ Further, every triangle has a unique Steiner circumellipse, which passes through the triangle's vertices and has its center at the triangle's centroid. Of all ellipses going through the triangle's vertices, it has the smallest area.
600
+
601
+ The Kiepert hyperbola is the unique conic which passes through the triangle's three vertices, its centroid, and its circumcenter.
602
+
603
+ Of all triangles contained in a given convex polygon, there exists a triangle with maximal area whose vertices are all vertices of the given polygon.[36]
604
+
605
+ One way to identify locations of points in (or outside) a triangle is to place the triangle in an arbitrary location and orientation in the Cartesian plane, and to use Cartesian coordinates. While convenient for many purposes, this approach has the disadvantage of all points' coordinate values being dependent on the arbitrary placement in the plane.
606
+
607
+ Two systems avoid that feature, so that the coordinates of a point are not affected by moving the triangle, rotating it, or reflecting it as in a mirror, any of which give a congruent triangle, or even by rescaling it to give a similar triangle:
608
+
609
+ A non-planar triangle is a triangle which is not contained in a (flat) plane. Some examples of non-planar triangles in non-Euclidean geometries are spherical triangles in spherical geometry and hyperbolic triangles in hyperbolic geometry.
610
+
611
+ While the measures of the internal angles in planar triangles always sum to 180°, a hyperbolic triangle has measures of angles that sum to less than 180°, and a spherical triangle has measures of angles that sum to more than 180°. A hyperbolic triangle can be obtained by drawing on a negatively curved surface, such as a saddle surface, and a spherical triangle can be obtained by drawing on a positively curved surface such as a sphere. Thus, if one draws a giant triangle on the surface of the Earth, one will find that the sum of the measures of its angles is greater than 180°; in fact it will be between 180° and 540°.[37] In particular it is possible to draw a triangle on a sphere such that the measure of each of its internal angles is equal to 90°, adding up to a total of 270°.
612
+
613
+ Specifically, on a sphere the sum of the angles of a triangle is
614
+
615
+ where f is the fraction of the sphere's area which is enclosed by the triangle. For example, suppose that we draw a triangle on the Earth's surface with vertices at the North Pole, at a point on the equator at 0° longitude, and a point on the equator at 90° West longitude. The great circle line between the latter two points is the equator, and the great circle line between either of those points and the North Pole is a line of longitude; so there are right angles at the two points on the equator. Moreover, the angle at the North Pole is also 90° because the other two vertices differ by 90° of longitude. So the sum of the angles in this triangle is 90° + 90° + 90° = 270°. The triangle encloses 1/4 of the northern hemisphere (90°/360° as viewed from the North Pole) and therefore 1/8 of the Earth's surface, so in the formula f = 1/8; thus the formula correctly gives the sum of the triangle's angles as 270°.
616
+
617
+ From the above angle sum formula we can also see that the Earth's surface is locally flat: If we draw an arbitrarily small triangle in the neighborhood of one point on the Earth's surface, the fraction f of the Earth's surface which is enclosed by the triangle will be arbitrarily close to zero. In this case the angle sum formula simplifies to 180°, which we know is what Euclidean geometry tells us for triangles on a flat surface.
618
+
619
+ Rectangles have been the most popular and common geometric form for buildings since the shape is easy to stack and organize; as a standard, it is easy to design furniture and fixtures to fit inside rectangularly shaped buildings. But triangles, while more difficult to use conceptually, provide a great deal of strength. As computer technology helps architects design creative new buildings, triangular shapes are becoming increasingly prevalent as parts of buildings and as the primary shape for some types of skyscrapers as well as building materials. In Tokyo in 1989, architects had wondered whether it was possible to build a 500-story tower to provide affordable office space for this densely packed city, but with the danger to buildings from earthquakes, architects considered that a triangular shape would be necessary if such a building were to be built.[38]
620
+
621
+ In New York City, as Broadway crisscrosses major avenues, the resulting blocks are cut like triangles, and buildings have been built on these shapes; one such building is the triangularly shaped Flatiron Building which real estate people admit has a "warren of awkward spaces that do not easily accommodate modern office furniture" but that has not prevented the structure from becoming a landmark icon.[39] Designers have made houses in Norway using triangular themes.[40] Triangle shapes have appeared in churches[41] as well as public buildings including colleges[42] as well as supports for innovative home designs.[43]
622
+
623
+ Triangles are sturdy; while a rectangle can collapse into a parallelogram from pressure to one of its points, triangles have a natural strength which supports structures against lateral pressures. A triangle will not change shape unless its sides are bent or extended or broken or if its joints break; in essence, each of the three sides supports the other two. A rectangle, in contrast, is more dependent on the strength of its joints in a structural sense. Some innovative designers have proposed making bricks not out of rectangles, but with triangular shapes which can be combined in three dimensions.[44] It is likely that triangles will be used increasingly in new ways as architecture increases in complexity. It is important to remember that triangles are strong in terms of rigidity, but while packed in a tessellating arrangement triangles are not as strong as hexagons under compression (hence the prevalence of hexagonal forms in nature). Tessellated triangles still maintain superior strength for cantilevering however, and this is the basis for one of the strongest man made structures, the tetrahedral truss.
624
+
625
+
626
+
627
+
628
+
en/5788.html.txt ADDED
@@ -0,0 +1,628 @@
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
1
+
2
+
3
+ A triangle is a polygon with three edges and three vertices. It is one of the basic shapes in geometry. A triangle with vertices A, B, and C is denoted
4
+
5
+
6
+
7
+ △
8
+ A
9
+ B
10
+ C
11
+
12
+
13
+ {\displaystyle \triangle ABC}
14
+
15
+ .
16
+
17
+ In Euclidean geometry any three points, when non-collinear, determine a unique triangle and simultaneously, a unique plane (i.e. a two-dimensional Euclidean space). In other words, there is only one plane that contains that triangle, and every triangle is contained in some plane. If the entire geometry is only the Euclidean plane, there is only one plane and all triangles are contained in it; however, in higher-dimensional Euclidean spaces, this is no longer true. This article is about triangles in Euclidean geometry, and in particular, the Euclidean plane, except where otherwise noted.
18
+
19
+ Triangles can be classified according to the lengths of their sides:
20
+
21
+ Hatch marks, also called tick marks, are used in diagrams of triangles and other geometric figures to identify sides of equal lengths. A side can be marked with a pattern of "ticks", short line segments in the form of tally marks; two sides have equal lengths if they are both marked with the same pattern. In a triangle, the pattern is usually no more than 3 ticks. An equilateral triangle has the same pattern on all 3 sides, an isosceles triangle has the same pattern on just 2 sides, and a scalene triangle has different patterns on all sides since no sides are equal. Similarly, patterns of 1, 2, or 3 concentric arcs inside the angles are used to indicate equal angles. An equilateral triangle has the same pattern on all 3 angles, an isosceles triangle has the same pattern on just 2 angles, and a scalene triangle has different patterns on all angles since no angles are equal.
22
+
23
+ Triangles can also be classified according to their internal angles, measured here in degrees.
24
+
25
+ A triangle that has two angles with the same measure also has two sides with the same length, and therefore it is an isosceles triangle. It follows that in a triangle where all angles have the same measure, all three sides have the same length, and such a triangle is therefore equilateral.
26
+
27
+ Triangles are assumed to be two-dimensional plane figures, unless the context provides otherwise (see Non-planar triangles, below). In rigorous treatments, a triangle is therefore called a 2-simplex (see also Polytope). Elementary facts about triangles were presented by Euclid in books 1–4 of his Elements, around 300 BC.
28
+
29
+ The sum of the measures of the interior angles of a triangle in Euclidean space is always 180 degrees.[5] This fact is equivalent to Euclid's parallel postulate. This allows determination of the measure of the third angle of any triangle given the measure of two angles. An exterior angle of a triangle is an angle that is a linear pair (and hence supplementary) to an interior angle. The measure of an exterior angle of a triangle is equal to the sum of the measures of the two interior angles that are not adjacent to it; this is the exterior angle theorem. The sum of the measures of the three exterior angles (one for each vertex) of any triangle is 360 degrees.[note 2]
30
+
31
+ Two triangles are said to be similar if every angle of one triangle has the same measure as the corresponding angle in the other triangle. The corresponding sides of similar triangles have lengths that are in the same proportion, and this property is also sufficient to establish similarity.
32
+
33
+ Some basic theorems about similar triangles are:
34
+
35
+ Two triangles that are congruent have exactly the same size and shape:[note 4] all pairs of corresponding interior angles are equal in measure, and all pairs of corresponding sides have the same length. (This is a total of six equalities, but three are often sufficient to prove congruence.)
36
+
37
+ Some individually necessary and sufficient conditions for a pair of triangles to be congruent are:
38
+
39
+ Some individually sufficient conditions are:
40
+
41
+ An important condition is:
42
+
43
+ Using right triangles and the concept of similarity, the trigonometric functions sine and cosine can be defined. These are functions of an angle which are investigated in trigonometry.
44
+
45
+ A central theorem is the Pythagorean theorem, which states in any right triangle, the square of the length of the hypotenuse equals the sum of the squares of the lengths of the two other sides. If the hypotenuse has length c, and the legs have lengths a and b, then the theorem states that
46
+
47
+ The converse is true: if the lengths of the sides of a triangle satisfy the above equation, then the triangle has a right angle opposite side c.
48
+
49
+ Some other facts about right triangles:
50
+
51
+ For all triangles, angles and sides are related by the law of cosines and law of sines (also called the cosine rule and sine rule).
52
+
53
+ The triangle inequality states that the sum of the lengths of any two sides of a triangle must be greater than or equal to the length of the third side. That sum can equal the length of the third side only in the case of a degenerate triangle, one with collinear vertices. It is not possible for that sum to be less than the length of the third side. A triangle with three given positive side lengths exists if and only if those side lengths satisfy the triangle inequality.
54
+
55
+ Three given angles form a non-degenerate triangle (and indeed an infinitude of them) if and only if both of these conditions hold: (a) each of the angles is positive, and (b) the angles sum to 180°. If degenerate triangles are permitted, angles of 0° are permitted.
56
+
57
+ Three positive angles α, β, and γ, each of them less than 180°, are the angles of a triangle if and only if any one of the following conditions holds:
58
+
59
+ the last equality applying only if none of the angles is 90° (so the tangent function's value is always finite).
60
+
61
+ There are thousands of different constructions that find a special point associated with (and often inside) a triangle, satisfying some unique property: see the article Encyclopedia of Triangle Centers for a catalogue of them. Often they are constructed by finding three lines associated in a symmetrical way with the three sides (or vertices) and then proving that the three lines meet in a single point: an important tool for proving the existence of these is Ceva's theorem, which gives a criterion for determining when three such lines are concurrent. Similarly, lines associated with a triangle are often constructed by proving that three symmetrically constructed points are collinear: here Menelaus' theorem gives a useful general criterion. In this section just a few of the most commonly encountered constructions are explained.
62
+
63
+ A perpendicular bisector of a side of a triangle is a straight line passing through the midpoint of the side and being perpendicular to it, i.e. forming a right angle with it. The three perpendicular bisectors meet in a single point, the triangle's circumcenter, usually denoted by O; this point is the center of the circumcircle, the circle passing through all three vertices. The diameter of this circle, called the circumdiameter, can be found from the law of sines stated above. The circumcircle's radius is called the circumradius.
64
+
65
+ Thales' theorem implies that if the circumcenter is located on a side of the triangle, then the opposite angle is a right one. If the circumcenter is located inside the triangle, then the triangle is acute; if the circumcenter is located outside the triangle, then the triangle is obtuse.
66
+
67
+ An altitude of a triangle is a straight line through a vertex and perpendicular to (i.e. forming a right angle with) the opposite side. This opposite side is called the base of the altitude, and the point where the altitude intersects the base (or its extension) is called the foot of the altitude. The length of the altitude is the distance between the base and the vertex. The three altitudes intersect in a single point, called the orthocenter of the triangle, usually denoted by H. The orthocenter lies inside the triangle if and only if the triangle is acute.
68
+
69
+ An angle bisector of a triangle is a straight line through a vertex which cuts the corresponding angle in half. The three angle bisectors intersect in a single point, the incenter, usually denoted by I, the center of the triangle's incircle. The incircle is the circle which lies inside the triangle and touches all three sides. Its radius is called the inradius. There are three other important circles, the excircles; they lie outside the triangle and touch one side as well as the extensions of the other two. The centers of the in- and excircles form an orthocentric system.
70
+
71
+ A median of a triangle is a straight line through a vertex and the midpoint of the opposite side, and divides the triangle into two equal areas. The three medians intersect in a single point, the triangle's centroid or geometric barycenter, usually denoted by G. The centroid of a rigid triangular object (cut out of a thin sheet of uniform density) is also its center of mass: the object can be balanced on its centroid in a uniform gravitational field. The centroid cuts every median in the ratio 2:1, i.e. the distance between a vertex and the centroid is twice the distance between the centroid and the midpoint of the opposite side.
72
+
73
+ The midpoints of the three sides and the feet of the three altitudes all lie on a single circle, the triangle's nine-point circle. The remaining three points for which it is named are the midpoints of the portion of altitude between the vertices and the orthocenter. The radius of the nine-point circle is half that of the circumcircle. It touches the incircle (at the Feuerbach point) and the three excircles.
74
+
75
+ The orthocenter (blue point), center of the nine-point circle (red), centroid (orange), and circumcenter (green) all lie on a single line, known as Euler's line (red line). The center of the nine-point circle lies at the midpoint between the orthocenter and the circumcenter, and the distance between the centroid and the circumcenter is half that between the centroid and the orthocenter.
76
+
77
+ The center of the incircle is not in general located on Euler's line.
78
+
79
+ If one reflects a median in the angle bisector that passes through the same vertex, one obtains a symmedian. The three symmedians intersect in a single point, the symmedian point of the triangle.
80
+
81
+ There are various standard methods for calculating the length of a side or the measure of an angle. Certain methods are suited to calculating values in a right-angled triangle; more complex methods may be required in other situations.
82
+
83
+ In right triangles, the trigonometric ratios of sine, cosine and tangent can be used to find unknown angles and the lengths of unknown sides. The sides of the triangle are known as follows:
84
+
85
+ The sine of an angle is the ratio of the length of the opposite side to the length of the hypotenuse. In our case
86
+
87
+ This ratio does not depend on the particular right triangle chosen, as long as it contains the angle A, since all those triangles are similar.
88
+
89
+ The cosine of an angle is the ratio of the length of the adjacent side to the length of the hypotenuse. In our case
90
+
91
+ The tangent of an angle is the ratio of the length of the opposite side to the length of the adjacent side. In our case
92
+
93
+ The acronym "SOH-CAH-TOA" is a useful mnemonic for these ratios.
94
+
95
+ The inverse trigonometric functions can be used to calculate the internal angles for a right angled triangle with the length of any two sides.
96
+
97
+ Arcsin can be used to calculate an angle from the length of the opposite side and the length of the hypotenuse.
98
+
99
+ Arccos can be used to calculate an angle from the length of the adjacent side and the length of the hypotenuse.
100
+
101
+ Arctan can be used to calculate an angle from the length of the opposite side and the length of the adjacent side.
102
+
103
+ In introductory geometry and trigonometry courses, the notation sin−1, cos−1, etc., are often used in place of arcsin, arccos, etc. However, the arcsin, arccos, etc., notation is standard in higher mathematics where trigonometric functions are commonly raised to powers, as this avoids confusion between multiplicative inverse and compositional inverse.
104
+
105
+ The law of sines, or sine rule,[8] states that the ratio of the length of a side to the sine of its corresponding opposite angle is constant, that is
106
+
107
+ This ratio is equal to the diameter of the circumscribed circle of the given triangle. Another interpretation of this theorem is that every triangle with angles α, β and γ is similar to a triangle with side lengths equal to sin α, sin β and sin γ. This triangle can be constructed by first constructing a circle of diameter 1, and inscribing in it two of the angles of the triangle. The length of the sides of that triangle will be sin α, sin β and sin γ. The side whose length is sin α is opposite to the angle whose measure is α, etc.
108
+
109
+ The law of cosines, or cosine rule, connects the length of an unknown side of a triangle to the length of the other sides and the angle opposite to the unknown side.[8] As per the law:
110
+
111
+ For a triangle with length of sides a, b, c and angles of α, β, γ respectively, given two known lengths of a triangle a and b, and the angle between the two known sides γ (or the angle opposite to the unknown side c), to calculate the third side c, the following formula can be used:
112
+
113
+ If the lengths of all three sides of any triangle are known the three angles can be calculated:
114
+
115
+ The law of tangents, or tangent rule, can be used to find a side or an angle when two sides and an angle or two angles and a side are known. It states that:[9]
116
+
117
+ "Solution of triangles" is the main trigonometric problem: to find missing characteristics of a triangle (three angles, the lengths of the three sides etc.) when at least three of these characteristics are given. The triangle can be located on a plane or on a sphere. This problem often occurs in various trigonometric applications, such as geodesy, astronomy, construction, navigation etc.
118
+
119
+ Calculating the area T of a triangle is an elementary problem encountered often in many different situations. The best known and simplest formula is:
120
+
121
+ where b is the length of the base of the triangle, and h is the height or altitude of the triangle. The term "base" denotes any side, and "height" denotes the length of a perpendicular from the vertex opposite the base onto the line containing the base. In 499 CE Aryabhata, used this illustrated method in the Aryabhatiya (section 2.6).[10]
122
+
123
+ Although simple, this formula is only useful if the height can be readily found, which is not always the case. For example, the surveyor of a triangular field might find it relatively easy to measure the length of each side, but relatively difficult to construct a 'height'. Various methods may be used in practice, depending on what is known about the triangle. The following is a selection of frequently used formulae for the area of a triangle.[11]
124
+
125
+ The height of a triangle can be found through the application of trigonometry.
126
+
127
+ Knowing SAS: Using the labels in the image on the right, the altitude is h = a sin
128
+
129
+
130
+
131
+ γ
132
+
133
+
134
+ {\displaystyle \gamma }
135
+
136
+ . Substituting this in the formula
137
+
138
+
139
+
140
+ T
141
+ =
142
+
143
+
144
+ 1
145
+ 2
146
+
147
+
148
+ b
149
+ h
150
+
151
+
152
+ {\displaystyle T={\frac {1}{2}}bh}
153
+
154
+ derived above, the area of the triangle can be expressed as:
155
+
156
+ (where α is the interior angle at A, β is the interior angle at B,
157
+
158
+
159
+
160
+ γ
161
+
162
+
163
+ {\displaystyle \gamma }
164
+
165
+ is the interior angle at C and c is the line AB).
166
+
167
+ Furthermore, since sin α = sin (π − α) = sin (β +
168
+
169
+
170
+
171
+ γ
172
+
173
+
174
+ {\displaystyle \gamma }
175
+
176
+ ), and similarly for the other two angles:
177
+
178
+ Knowing AAS:
179
+
180
+ and analogously if the known side is a or c.
181
+
182
+ Knowing ASA:[12]
183
+
184
+ and analogously if the known side is b or c.
185
+
186
+ The shape of the triangle is determined by the lengths of the sides. Therefore, the area can also be derived from the lengths of the sides. By Heron's formula:
187
+
188
+ where
189
+
190
+
191
+
192
+ s
193
+ =
194
+
195
+
196
+
197
+
198
+ a
199
+ +
200
+ b
201
+ +
202
+ c
203
+
204
+ 2
205
+
206
+
207
+
208
+
209
+
210
+ {\displaystyle s={\tfrac {a+b+c}{2}}}
211
+
212
+ is the semiperimeter, or half of the triangle's perimeter.
213
+
214
+ Three other equivalent ways of writing Heron's formula are
215
+
216
+ The area of a parallelogram embedded in a three-dimensional Euclidean space can be calculated using vectors. Let vectors AB and AC point respectively from A to B and from A to C. The area of parallelogram ABDC is then
217
+
218
+ which is the magnitude of the cross product of vectors AB and AC. The area of triangle ABC is half of this,
219
+
220
+ The area of triangle ABC can also be expressed in terms of dot products as follows:
221
+
222
+ In two-dimensional Euclidean space, expressing vector AB as a free vector in Cartesian space equal to (x1,y1) and AC as (x2,y2), this can be rewritten as:
223
+
224
+ If vertex A is located at the origin (0, 0) of a Cartesian coordinate system and the coordinates of the other two vertices are given by B = (xB, yB) and C = (xC, yC), then the area can be computed as ​1⁄2 times the absolute value of the determinant
225
+
226
+ For three general vertices, the equation is:
227
+
228
+ which can be written as
229
+
230
+ If the points are labeled sequentially in the counterclockwise direction, the above determinant expressions are positive and the absolute value signs can be omitted.[13] The above formula is known as the shoelace formula or the surveyor's formula.
231
+
232
+ If we locate the vertices in the complex plane and denote them in counterclockwise sequence as a = xA + yAi, b = xB + yBi, and c = xC + yCi, and denote their complex conjugates as
233
+
234
+
235
+
236
+
237
+
238
+
239
+ a
240
+ ¯
241
+
242
+
243
+
244
+
245
+
246
+ {\displaystyle {\bar {a}}}
247
+
248
+ ,
249
+
250
+
251
+
252
+
253
+
254
+
255
+ b
256
+ ¯
257
+
258
+
259
+
260
+
261
+
262
+ {\displaystyle {\bar {b}}}
263
+
264
+ , and
265
+
266
+
267
+
268
+
269
+
270
+
271
+ c
272
+ ¯
273
+
274
+
275
+
276
+
277
+
278
+ {\displaystyle {\bar {c}}}
279
+
280
+ , then the formula
281
+
282
+ is equivalent to the shoelace formula.
283
+
284
+ In three dimensions, the area of a general triangle A = (xA, yA, zA), B = (xB, yB, zB) and C = (xC, yC, zC) is the Pythagorean sum of the areas of the respective projections on the three principal planes (i.e. x = 0, y = 0 and z = 0):
285
+
286
+ The area within any closed curve, such as a triangle, is given by the line integral around the curve of the algebraic or signed distance of a point on the curve from an arbitrary oriented straight line L. Points to the right of L as oriented are taken to be at negative distance from L, while the weight for the integral is taken to be the component of arc length parallel to L rather than arc length itself.
287
+
288
+ This method is well suited to computation of the area of an arbitrary polygon. Taking L to be the x-axis, the line integral between consecutive vertices (xi,yi) and (xi+1,yi+1) is given by the base times the mean height, namely (xi+1 − xi)(yi + yi+1)/2. The sign of the area is an overall indicator of the direction of traversal, with negative area indicating counterclockwise traversal. The area of a triangle then falls out as the case of a polygon with three sides.
289
+
290
+ While the line integral method has in common with other coordinate-based methods the arbitrary choice of a coordinate system, unlike the others it makes no arbitrary choice of vertex of the triangle as origin or of side as base. Furthermore, the choice of coordinate system defined by L commits to only two degrees of freedom rather than the usual three, since the weight is a local distance (e.g. xi+1 − xi in the above) whence the method does not require choosing an axis normal to L.
291
+
292
+ When working in polar coordinates it is not necessary to convert to Cartesian coordinates to use line integration, since the line integral between consecutive vertices (ri,θi) and (ri+1,θi+1) of a polygon is given directly by riri+1sin(θi+1 − θi)/2. This is valid for all values of θ, with some decrease in numerical accuracy when |θ| is many orders of magnitude greater than π. With this formulation negative area indicates clockwise traversal, which should be kept in mind when mixing polar and cartesian coordinates. Just as the choice of y-axis (x = 0) is immaterial for line integration in cartesian coordinates, so is the choice of zero heading (θ = 0) immaterial here.
293
+
294
+ Three formulas have the same structure as Heron's formula but are expressed in terms of different variables. First, denoting the medians from sides a, b, and c respectively as ma, mb, and mc and their semi-sum (ma + mb + mc)/2 as σ, we have[14]
295
+
296
+ Next, denoting the altitudes from sides a, b, and c respectively as ha, hb, and hc, and denoting the semi-sum of the reciprocals of the altitudes as
297
+
298
+
299
+
300
+ H
301
+ =
302
+ (
303
+
304
+ h
305
+
306
+ a
307
+
308
+
309
+ −
310
+ 1
311
+
312
+
313
+ +
314
+
315
+ h
316
+
317
+ b
318
+
319
+
320
+ −
321
+ 1
322
+
323
+
324
+ +
325
+
326
+ h
327
+
328
+ c
329
+
330
+
331
+ −
332
+ 1
333
+
334
+
335
+ )
336
+
337
+ /
338
+
339
+ 2
340
+
341
+
342
+ {\displaystyle H=(h_{a}^{-1}+h_{b}^{-1}+h_{c}^{-1})/2}
343
+
344
+ we have[15]
345
+
346
+ And denoting the semi-sum of the angles' sines as S = [(sin α) + (sin β) + (sin γ)]/2, we have[16]
347
+
348
+ where D is the diameter of the circumcircle:
349
+
350
+
351
+
352
+ D
353
+ =
354
+
355
+
356
+
357
+ a
358
+
359
+ sin
360
+ ⁡
361
+ α
362
+
363
+
364
+
365
+
366
+ =
367
+
368
+
369
+
370
+ b
371
+
372
+ sin
373
+ ⁡
374
+ β
375
+
376
+
377
+
378
+
379
+ =
380
+
381
+
382
+
383
+ c
384
+
385
+ sin
386
+ ⁡
387
+ γ
388
+
389
+
390
+
391
+
392
+ .
393
+
394
+
395
+ {\displaystyle D={\tfrac {a}{\sin \alpha }}={\tfrac {b}{\sin \beta }}={\tfrac {c}{\sin \gamma }}.}
396
+
397
+ See Pick's theorem for a technique for finding the area of any arbitrary lattice polygon (one drawn on a grid with vertically and horizontally adjacent lattice points at equal distances, and with vertices on lattice points).
398
+
399
+ The theorem states:
400
+
401
+ where
402
+
403
+
404
+
405
+ I
406
+
407
+
408
+ {\displaystyle I}
409
+
410
+ is the number of internal lattice points and B is the number of lattice points lying on the border of the polygon.
411
+
412
+ Numerous other area formulas exist, such as
413
+
414
+ where r is the inradius, and s is the semiperimeter (in fact, this formula holds for all tangential polygons), and[17]:Lemma 2
415
+
416
+ where
417
+
418
+
419
+
420
+
421
+ r
422
+
423
+ a
424
+
425
+
426
+ ,
427
+
428
+
429
+ r
430
+
431
+ b
432
+
433
+
434
+ ,
435
+
436
+
437
+ r
438
+
439
+ c
440
+
441
+
442
+
443
+
444
+ {\displaystyle r_{a},\,r_{b},\,r_{c}}
445
+
446
+ are the radii of the excircles tangent to sides a, b, c respectively.
447
+
448
+ We also have
449
+
450
+ and[18]
451
+
452
+ for circumdiameter D; and[19]
453
+
454
+ for angle α ≠ 90°.
455
+
456
+ The area can also be expressed as[20]
457
+
458
+ In 1885, Baker[21] gave a collection of over a hundred distinct area formulas for the triangle. These include:
459
+
460
+ for circumradius (radius of the circumcircle) R, and
461
+
462
+ The area T of any triangle with perimeter p satisfies
463
+
464
+ with equality holding if and only if the triangle is equilateral.[22][23]:657
465
+
466
+ Other upper bounds on the area T are given by[24]:p.290
467
+
468
+ and
469
+
470
+ both again holding if and only if the triangle is equilateral.
471
+
472
+ There are infinitely many lines that bisect the area of a triangle.[25] Three of them are the medians, which are the only area bisectors that go through the centroid. Three other area bisectors are parallel to the triangle's sides.
473
+
474
+ Any line through a triangle that splits both the triangle's area and its perimeter in half goes through the triangle's incenter. There can be one, two, or three of these for any given triangle.
475
+
476
+ The formulas in this section are true for all Euclidean triangles.
477
+
478
+ The medians and the sides are related by[26]:p.70
479
+
480
+ and
481
+
482
+ and equivalently for mb and mc.
483
+
484
+ For angle A opposite side a, the length of the internal angle bisector is given by[27]
485
+
486
+ for semiperimeter s, where the bisector length is measured from the vertex to where it meets the opposite side.
487
+
488
+ The interior perpendicular bisectors are given by
489
+
490
+ where the sides are
491
+
492
+
493
+
494
+ a
495
+ ≥
496
+ b
497
+ ≥
498
+ c
499
+
500
+
501
+ {\displaystyle a\geq b\geq c}
502
+
503
+ and the area is
504
+
505
+
506
+
507
+ T
508
+ .
509
+
510
+
511
+ {\displaystyle T.}
512
+
513
+ [28]:Thm 2
514
+
515
+ The altitude from, for example, the side of length a is
516
+
517
+ The following formulas involve the circumradius R and the inradius r:
518
+
519
+ where ha etc. are the altitudes to the subscripted sides;[26]:p.79
520
+
521
+ and
522
+
523
+ The product of two sides of a triangle equals the altitude to the third side times the diameter D of the circumcircle:[26]:p.64
524
+
525
+ Suppose two adjacent but non-overlapping triangles share the same side of length f and share the same circumcircle, so that the side of length f is a chord of the circumcircle and the triangles have side lengths (a, b, f) and (c, d, f), with the two triangles together forming a cyclic quadrilateral with side lengths in sequence (a, b, c, d). Then[29]:84
526
+
527
+ Let G be the centroid of a triangle with vertices A, B, and C, and let P be any interior point. Then the distances between the points are related by[29]:174
528
+
529
+ The sum of the squares of the triangle's sides equals three times the sum of the squared distances of the centroid from the vertices:
530
+
531
+ Let qa, qb, and qc be the distances from the centroid to the sides of lengths a, b, and c. Then[29]:173
532
+
533
+ and
534
+
535
+ for area T.
536
+
537
+ Carnot's theorem states that the sum of the distances from the circumcenter to the three sides equals the sum of the circumradius and the inradius.[26]:p.83 Here a segment's length is considered to be negative if and only if the segment lies entirely outside the triangle. This method is especially useful for deducing the properties of more abstract forms of triangles, such as the ones induced by Lie algebras, that otherwise have the same properties as usual triangles.
538
+
539
+ Euler's theorem states that the distance d between the circumcenter and the incenter is given by[26]:p.85
540
+
541
+ or equivalently
542
+
543
+ where R is the circumradius and r is the inradius. Thus for all triangles R ≥ 2r, with equality holding for equilateral triangles.
544
+
545
+ If we denote that the orthocenter divides one altitude into segments of lengths u and v, another altitude into segment lengths w and x, and the third altitude into segment lengths y and z, then uv = wx = yz.[26]:p.94
546
+
547
+ The distance from a side to the circumcenter equals half the distance from the opposite vertex to the orthocenter.[26]:p.99
548
+
549
+ The sum of the squares of the distances from the vertices to the orthocenter H plus the sum of the squares of the sides equals twelve times the square of the circumradius:[26]:p.102
550
+
551
+ In addition to the law of sines, the law of cosines, the law of tangents, and the trigonometric existence conditions given earlier, for any triangle
552
+
553
+ Morley's trisector theorem states that in any triangle, the three points of intersection of the adjacent angle trisectors form an equilateral triangle, called the Morley triangle.
554
+
555
+ As discussed above, every triangle has a unique inscribed circle (incircle) that is interior to the triangle and tangent to all three sides.
556
+
557
+ Every triangle has a unique Steiner inellipse which is interior to the triangle and tangent at the midpoints of the sides. Marden's theorem shows how to find the foci of this ellipse.[31] This ellipse has the greatest area of any ellipse tangent to all three sides of the triangle.
558
+
559
+ The Mandart inellipse of a triangle is the ellipse inscribed within the triangle tangent to its sides at the contact points of its excircles.
560
+
561
+ For any ellipse inscribed in a triangle ABC, let the foci be P and Q. Then[32]
562
+
563
+ Every convex polygon with area T can be inscribed in a triangle of area at most equal to 2T. Equality holds (exclusively) for a parallelogram.[33]
564
+
565
+ The Lemoine hexagon is a cyclic hexagon with vertices given by the six intersections of the sides of a triangle with the three lines that are parallel to the sides and that pass through its symmedian point. In either its simple form or its self-intersecting form, the Lemoine hexagon is interior to the triangle with two vertices on each side of the triangle.
566
+
567
+ Every acute triangle has three inscribed squares (squares in its interior such that all four of a square's vertices lie on a side of the triangle, so two of them lie on the same side and hence one side of the square coincides with part of a side of the triangle). In a right triangle two of the squares coincide and have a vertex at the triangle's right angle, so a right triangle has only two distinct inscribed squares. An obtuse triangle has only one inscribed square, with a side coinciding with part of the triangle's longest side. Within a given triangle, a longer common side is associated with a smaller inscribed square. If an inscribed square has side of length qa and the triangle has a side of length a, part of which side coincides with a side of the square, then qa, a, the altitude ha from the side a, and the triangle's area T are related according to[34][35]
568
+
569
+ The largest possible ratio of the area of the inscribed square to the area of the triangle is 1/2, which occurs when a2 = 2T, q = a/2, and the altitude of the triangle from the base of length a is equal to a. The smallest possible ratio of the side of one inscribed square to the side of another in the same non-obtuse triangle is
570
+
571
+
572
+
573
+ 2
574
+
575
+
576
+ 2
577
+
578
+
579
+
580
+ /
581
+
582
+ 3
583
+ =
584
+ 0.94....
585
+
586
+
587
+ {\displaystyle 2{\sqrt {2}}/3=0.94....}
588
+
589
+ [35] Both of these extreme cases occur for the isosceles right triangle.
590
+
591
+ From an interior point in a reference triangle, the nearest points on the three sides serve as the vertices of the pedal triangle of that point. If the interior point is the circumcenter of the reference triangle, the vertices of the pedal triangle are the midpoints of the reference triangle's sides, and so the pedal triangle is called the midpoint triangle or medial triangle. The midpoint triangle subdivides the reference triangle into four congruent triangles which are similar to the reference triangle.
592
+
593
+ The Gergonne triangle or intouch triangle of a reference triangle has its vertices at the three points of tangency of the reference triangle's sides with its incircle. The extouch triangle of a reference triangle has its vertices at the points of tangency of the reference triangle's excircles with its sides (not extended).
594
+
595
+ The tangential triangle of a reference triangle (other than a right triangle) is the triangle whose sides are on the tangent lines to the reference triangle's circumcircle at its vertices.
596
+
597
+ As mentioned above, every triangle has a unique circumcircle, a circle passing through all three vertices, whose center is the intersection of the perpendicular bisectors of the triangle's sides.
598
+
599
+ Further, every triangle has a unique Steiner circumellipse, which passes through the triangle's vertices and has its center at the triangle's centroid. Of all ellipses going through the triangle's vertices, it has the smallest area.
600
+
601
+ The Kiepert hyperbola is the unique conic which passes through the triangle's three vertices, its centroid, and its circumcenter.
602
+
603
+ Of all triangles contained in a given convex polygon, there exists a triangle with maximal area whose vertices are all vertices of the given polygon.[36]
604
+
605
+ One way to identify locations of points in (or outside) a triangle is to place the triangle in an arbitrary location and orientation in the Cartesian plane, and to use Cartesian coordinates. While convenient for many purposes, this approach has the disadvantage of all points' coordinate values being dependent on the arbitrary placement in the plane.
606
+
607
+ Two systems avoid that feature, so that the coordinates of a point are not affected by moving the triangle, rotating it, or reflecting it as in a mirror, any of which give a congruent triangle, or even by rescaling it to give a similar triangle:
608
+
609
+ A non-planar triangle is a triangle which is not contained in a (flat) plane. Some examples of non-planar triangles in non-Euclidean geometries are spherical triangles in spherical geometry and hyperbolic triangles in hyperbolic geometry.
610
+
611
+ While the measures of the internal angles in planar triangles always sum to 180°, a hyperbolic triangle has measures of angles that sum to less than 180°, and a spherical triangle has measures of angles that sum to more than 180°. A hyperbolic triangle can be obtained by drawing on a negatively curved surface, such as a saddle surface, and a spherical triangle can be obtained by drawing on a positively curved surface such as a sphere. Thus, if one draws a giant triangle on the surface of the Earth, one will find that the sum of the measures of its angles is greater than 180°; in fact it will be between 180° and 540°.[37] In particular it is possible to draw a triangle on a sphere such that the measure of each of its internal angles is equal to 90°, adding up to a total of 270°.
612
+
613
+ Specifically, on a sphere the sum of the angles of a triangle is
614
+
615
+ where f is the fraction of the sphere's area which is enclosed by the triangle. For example, suppose that we draw a triangle on the Earth's surface with vertices at the North Pole, at a point on the equator at 0° longitude, and a point on the equator at 90° West longitude. The great circle line between the latter two points is the equator, and the great circle line between either of those points and the North Pole is a line of longitude; so there are right angles at the two points on the equator. Moreover, the angle at the North Pole is also 90° because the other two vertices differ by 90° of longitude. So the sum of the angles in this triangle is 90° + 90° + 90° = 270°. The triangle encloses 1/4 of the northern hemisphere (90°/360° as viewed from the North Pole) and therefore 1/8 of the Earth's surface, so in the formula f = 1/8; thus the formula correctly gives the sum of the triangle's angles as 270°.
616
+
617
+ From the above angle sum formula we can also see that the Earth's surface is locally flat: If we draw an arbitrarily small triangle in the neighborhood of one point on the Earth's surface, the fraction f of the Earth's surface which is enclosed by the triangle will be arbitrarily close to zero. In this case the angle sum formula simplifies to 180°, which we know is what Euclidean geometry tells us for triangles on a flat surface.
618
+
619
+ Rectangles have been the most popular and common geometric form for buildings since the shape is easy to stack and organize; as a standard, it is easy to design furniture and fixtures to fit inside rectangularly shaped buildings. But triangles, while more difficult to use conceptually, provide a great deal of strength. As computer technology helps architects design creative new buildings, triangular shapes are becoming increasingly prevalent as parts of buildings and as the primary shape for some types of skyscrapers as well as building materials. In Tokyo in 1989, architects had wondered whether it was possible to build a 500-story tower to provide affordable office space for this densely packed city, but with the danger to buildings from earthquakes, architects considered that a triangular shape would be necessary if such a building were to be built.[38]
620
+
621
+ In New York City, as Broadway crisscrosses major avenues, the resulting blocks are cut like triangles, and buildings have been built on these shapes; one such building is the triangularly shaped Flatiron Building which real estate people admit has a "warren of awkward spaces that do not easily accommodate modern office furniture" but that has not prevented the structure from becoming a landmark icon.[39] Designers have made houses in Norway using triangular themes.[40] Triangle shapes have appeared in churches[41] as well as public buildings including colleges[42] as well as supports for innovative home designs.[43]
622
+
623
+ Triangles are sturdy; while a rectangle can collapse into a parallelogram from pressure to one of its points, triangles have a natural strength which supports structures against lateral pressures. A triangle will not change shape unless its sides are bent or extended or broken or if its joints break; in essence, each of the three sides supports the other two. A rectangle, in contrast, is more dependent on the strength of its joints in a structural sense. Some innovative designers have proposed making bricks not out of rectangles, but with triangular shapes which can be combined in three dimensions.[44] It is likely that triangles will be used increasingly in new ways as architecture increases in complexity. It is important to remember that triangles are strong in terms of rigidity, but while packed in a tessellating arrangement triangles are not as strong as hexagons under compression (hence the prevalence of hexagonal forms in nature). Tessellated triangles still maintain superior strength for cantilevering however, and this is the basis for one of the strongest man made structures, the tetrahedral truss.
624
+
625
+
626
+
627
+
628
+
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1
+
2
+
3
+ A triangle is a polygon with three edges and three vertices. It is one of the basic shapes in geometry. A triangle with vertices A, B, and C is denoted
4
+
5
+
6
+
7
+ △
8
+ A
9
+ B
10
+ C
11
+
12
+
13
+ {\displaystyle \triangle ABC}
14
+
15
+ .
16
+
17
+ In Euclidean geometry any three points, when non-collinear, determine a unique triangle and simultaneously, a unique plane (i.e. a two-dimensional Euclidean space). In other words, there is only one plane that contains that triangle, and every triangle is contained in some plane. If the entire geometry is only the Euclidean plane, there is only one plane and all triangles are contained in it; however, in higher-dimensional Euclidean spaces, this is no longer true. This article is about triangles in Euclidean geometry, and in particular, the Euclidean plane, except where otherwise noted.
18
+
19
+ Triangles can be classified according to the lengths of their sides:
20
+
21
+ Hatch marks, also called tick marks, are used in diagrams of triangles and other geometric figures to identify sides of equal lengths. A side can be marked with a pattern of "ticks", short line segments in the form of tally marks; two sides have equal lengths if they are both marked with the same pattern. In a triangle, the pattern is usually no more than 3 ticks. An equilateral triangle has the same pattern on all 3 sides, an isosceles triangle has the same pattern on just 2 sides, and a scalene triangle has different patterns on all sides since no sides are equal. Similarly, patterns of 1, 2, or 3 concentric arcs inside the angles are used to indicate equal angles. An equilateral triangle has the same pattern on all 3 angles, an isosceles triangle has the same pattern on just 2 angles, and a scalene triangle has different patterns on all angles since no angles are equal.
22
+
23
+ Triangles can also be classified according to their internal angles, measured here in degrees.
24
+
25
+ A triangle that has two angles with the same measure also has two sides with the same length, and therefore it is an isosceles triangle. It follows that in a triangle where all angles have the same measure, all three sides have the same length, and such a triangle is therefore equilateral.
26
+
27
+ Triangles are assumed to be two-dimensional plane figures, unless the context provides otherwise (see Non-planar triangles, below). In rigorous treatments, a triangle is therefore called a 2-simplex (see also Polytope). Elementary facts about triangles were presented by Euclid in books 1–4 of his Elements, around 300 BC.
28
+
29
+ The sum of the measures of the interior angles of a triangle in Euclidean space is always 180 degrees.[5] This fact is equivalent to Euclid's parallel postulate. This allows determination of the measure of the third angle of any triangle given the measure of two angles. An exterior angle of a triangle is an angle that is a linear pair (and hence supplementary) to an interior angle. The measure of an exterior angle of a triangle is equal to the sum of the measures of the two interior angles that are not adjacent to it; this is the exterior angle theorem. The sum of the measures of the three exterior angles (one for each vertex) of any triangle is 360 degrees.[note 2]
30
+
31
+ Two triangles are said to be similar if every angle of one triangle has the same measure as the corresponding angle in the other triangle. The corresponding sides of similar triangles have lengths that are in the same proportion, and this property is also sufficient to establish similarity.
32
+
33
+ Some basic theorems about similar triangles are:
34
+
35
+ Two triangles that are congruent have exactly the same size and shape:[note 4] all pairs of corresponding interior angles are equal in measure, and all pairs of corresponding sides have the same length. (This is a total of six equalities, but three are often sufficient to prove congruence.)
36
+
37
+ Some individually necessary and sufficient conditions for a pair of triangles to be congruent are:
38
+
39
+ Some individually sufficient conditions are:
40
+
41
+ An important condition is:
42
+
43
+ Using right triangles and the concept of similarity, the trigonometric functions sine and cosine can be defined. These are functions of an angle which are investigated in trigonometry.
44
+
45
+ A central theorem is the Pythagorean theorem, which states in any right triangle, the square of the length of the hypotenuse equals the sum of the squares of the lengths of the two other sides. If the hypotenuse has length c, and the legs have lengths a and b, then the theorem states that
46
+
47
+ The converse is true: if the lengths of the sides of a triangle satisfy the above equation, then the triangle has a right angle opposite side c.
48
+
49
+ Some other facts about right triangles:
50
+
51
+ For all triangles, angles and sides are related by the law of cosines and law of sines (also called the cosine rule and sine rule).
52
+
53
+ The triangle inequality states that the sum of the lengths of any two sides of a triangle must be greater than or equal to the length of the third side. That sum can equal the length of the third side only in the case of a degenerate triangle, one with collinear vertices. It is not possible for that sum to be less than the length of the third side. A triangle with three given positive side lengths exists if and only if those side lengths satisfy the triangle inequality.
54
+
55
+ Three given angles form a non-degenerate triangle (and indeed an infinitude of them) if and only if both of these conditions hold: (a) each of the angles is positive, and (b) the angles sum to 180°. If degenerate triangles are permitted, angles of 0° are permitted.
56
+
57
+ Three positive angles α, β, and γ, each of them less than 180°, are the angles of a triangle if and only if any one of the following conditions holds:
58
+
59
+ the last equality applying only if none of the angles is 90° (so the tangent function's value is always finite).
60
+
61
+ There are thousands of different constructions that find a special point associated with (and often inside) a triangle, satisfying some unique property: see the article Encyclopedia of Triangle Centers for a catalogue of them. Often they are constructed by finding three lines associated in a symmetrical way with the three sides (or vertices) and then proving that the three lines meet in a single point: an important tool for proving the existence of these is Ceva's theorem, which gives a criterion for determining when three such lines are concurrent. Similarly, lines associated with a triangle are often constructed by proving that three symmetrically constructed points are collinear: here Menelaus' theorem gives a useful general criterion. In this section just a few of the most commonly encountered constructions are explained.
62
+
63
+ A perpendicular bisector of a side of a triangle is a straight line passing through the midpoint of the side and being perpendicular to it, i.e. forming a right angle with it. The three perpendicular bisectors meet in a single point, the triangle's circumcenter, usually denoted by O; this point is the center of the circumcircle, the circle passing through all three vertices. The diameter of this circle, called the circumdiameter, can be found from the law of sines stated above. The circumcircle's radius is called the circumradius.
64
+
65
+ Thales' theorem implies that if the circumcenter is located on a side of the triangle, then the opposite angle is a right one. If the circumcenter is located inside the triangle, then the triangle is acute; if the circumcenter is located outside the triangle, then the triangle is obtuse.
66
+
67
+ An altitude of a triangle is a straight line through a vertex and perpendicular to (i.e. forming a right angle with) the opposite side. This opposite side is called the base of the altitude, and the point where the altitude intersects the base (or its extension) is called the foot of the altitude. The length of the altitude is the distance between the base and the vertex. The three altitudes intersect in a single point, called the orthocenter of the triangle, usually denoted by H. The orthocenter lies inside the triangle if and only if the triangle is acute.
68
+
69
+ An angle bisector of a triangle is a straight line through a vertex which cuts the corresponding angle in half. The three angle bisectors intersect in a single point, the incenter, usually denoted by I, the center of the triangle's incircle. The incircle is the circle which lies inside the triangle and touches all three sides. Its radius is called the inradius. There are three other important circles, the excircles; they lie outside the triangle and touch one side as well as the extensions of the other two. The centers of the in- and excircles form an orthocentric system.
70
+
71
+ A median of a triangle is a straight line through a vertex and the midpoint of the opposite side, and divides the triangle into two equal areas. The three medians intersect in a single point, the triangle's centroid or geometric barycenter, usually denoted by G. The centroid of a rigid triangular object (cut out of a thin sheet of uniform density) is also its center of mass: the object can be balanced on its centroid in a uniform gravitational field. The centroid cuts every median in the ratio 2:1, i.e. the distance between a vertex and the centroid is twice the distance between the centroid and the midpoint of the opposite side.
72
+
73
+ The midpoints of the three sides and the feet of the three altitudes all lie on a single circle, the triangle's nine-point circle. The remaining three points for which it is named are the midpoints of the portion of altitude between the vertices and the orthocenter. The radius of the nine-point circle is half that of the circumcircle. It touches the incircle (at the Feuerbach point) and the three excircles.
74
+
75
+ The orthocenter (blue point), center of the nine-point circle (red), centroid (orange), and circumcenter (green) all lie on a single line, known as Euler's line (red line). The center of the nine-point circle lies at the midpoint between the orthocenter and the circumcenter, and the distance between the centroid and the circumcenter is half that between the centroid and the orthocenter.
76
+
77
+ The center of the incircle is not in general located on Euler's line.
78
+
79
+ If one reflects a median in the angle bisector that passes through the same vertex, one obtains a symmedian. The three symmedians intersect in a single point, the symmedian point of the triangle.
80
+
81
+ There are various standard methods for calculating the length of a side or the measure of an angle. Certain methods are suited to calculating values in a right-angled triangle; more complex methods may be required in other situations.
82
+
83
+ In right triangles, the trigonometric ratios of sine, cosine and tangent can be used to find unknown angles and the lengths of unknown sides. The sides of the triangle are known as follows:
84
+
85
+ The sine of an angle is the ratio of the length of the opposite side to the length of the hypotenuse. In our case
86
+
87
+ This ratio does not depend on the particular right triangle chosen, as long as it contains the angle A, since all those triangles are similar.
88
+
89
+ The cosine of an angle is the ratio of the length of the adjacent side to the length of the hypotenuse. In our case
90
+
91
+ The tangent of an angle is the ratio of the length of the opposite side to the length of the adjacent side. In our case
92
+
93
+ The acronym "SOH-CAH-TOA" is a useful mnemonic for these ratios.
94
+
95
+ The inverse trigonometric functions can be used to calculate the internal angles for a right angled triangle with the length of any two sides.
96
+
97
+ Arcsin can be used to calculate an angle from the length of the opposite side and the length of the hypotenuse.
98
+
99
+ Arccos can be used to calculate an angle from the length of the adjacent side and the length of the hypotenuse.
100
+
101
+ Arctan can be used to calculate an angle from the length of the opposite side and the length of the adjacent side.
102
+
103
+ In introductory geometry and trigonometry courses, the notation sin−1, cos−1, etc., are often used in place of arcsin, arccos, etc. However, the arcsin, arccos, etc., notation is standard in higher mathematics where trigonometric functions are commonly raised to powers, as this avoids confusion between multiplicative inverse and compositional inverse.
104
+
105
+ The law of sines, or sine rule,[8] states that the ratio of the length of a side to the sine of its corresponding opposite angle is constant, that is
106
+
107
+ This ratio is equal to the diameter of the circumscribed circle of the given triangle. Another interpretation of this theorem is that every triangle with angles α, β and γ is similar to a triangle with side lengths equal to sin α, sin β and sin γ. This triangle can be constructed by first constructing a circle of diameter 1, and inscribing in it two of the angles of the triangle. The length of the sides of that triangle will be sin α, sin β and sin γ. The side whose length is sin α is opposite to the angle whose measure is α, etc.
108
+
109
+ The law of cosines, or cosine rule, connects the length of an unknown side of a triangle to the length of the other sides and the angle opposite to the unknown side.[8] As per the law:
110
+
111
+ For a triangle with length of sides a, b, c and angles of α, β, γ respectively, given two known lengths of a triangle a and b, and the angle between the two known sides γ (or the angle opposite to the unknown side c), to calculate the third side c, the following formula can be used:
112
+
113
+ If the lengths of all three sides of any triangle are known the three angles can be calculated:
114
+
115
+ The law of tangents, or tangent rule, can be used to find a side or an angle when two sides and an angle or two angles and a side are known. It states that:[9]
116
+
117
+ "Solution of triangles" is the main trigonometric problem: to find missing characteristics of a triangle (three angles, the lengths of the three sides etc.) when at least three of these characteristics are given. The triangle can be located on a plane or on a sphere. This problem often occurs in various trigonometric applications, such as geodesy, astronomy, construction, navigation etc.
118
+
119
+ Calculating the area T of a triangle is an elementary problem encountered often in many different situations. The best known and simplest formula is:
120
+
121
+ where b is the length of the base of the triangle, and h is the height or altitude of the triangle. The term "base" denotes any side, and "height" denotes the length of a perpendicular from the vertex opposite the base onto the line containing the base. In 499 CE Aryabhata, used this illustrated method in the Aryabhatiya (section 2.6).[10]
122
+
123
+ Although simple, this formula is only useful if the height can be readily found, which is not always the case. For example, the surveyor of a triangular field might find it relatively easy to measure the length of each side, but relatively difficult to construct a 'height'. Various methods may be used in practice, depending on what is known about the triangle. The following is a selection of frequently used formulae for the area of a triangle.[11]
124
+
125
+ The height of a triangle can be found through the application of trigonometry.
126
+
127
+ Knowing SAS: Using the labels in the image on the right, the altitude is h = a sin
128
+
129
+
130
+
131
+ γ
132
+
133
+
134
+ {\displaystyle \gamma }
135
+
136
+ . Substituting this in the formula
137
+
138
+
139
+
140
+ T
141
+ =
142
+
143
+
144
+ 1
145
+ 2
146
+
147
+
148
+ b
149
+ h
150
+
151
+
152
+ {\displaystyle T={\frac {1}{2}}bh}
153
+
154
+ derived above, the area of the triangle can be expressed as:
155
+
156
+ (where α is the interior angle at A, β is the interior angle at B,
157
+
158
+
159
+
160
+ γ
161
+
162
+
163
+ {\displaystyle \gamma }
164
+
165
+ is the interior angle at C and c is the line AB).
166
+
167
+ Furthermore, since sin α = sin (π − α) = sin (β +
168
+
169
+
170
+
171
+ γ
172
+
173
+
174
+ {\displaystyle \gamma }
175
+
176
+ ), and similarly for the other two angles:
177
+
178
+ Knowing AAS:
179
+
180
+ and analogously if the known side is a or c.
181
+
182
+ Knowing ASA:[12]
183
+
184
+ and analogously if the known side is b or c.
185
+
186
+ The shape of the triangle is determined by the lengths of the sides. Therefore, the area can also be derived from the lengths of the sides. By Heron's formula:
187
+
188
+ where
189
+
190
+
191
+
192
+ s
193
+ =
194
+
195
+
196
+
197
+
198
+ a
199
+ +
200
+ b
201
+ +
202
+ c
203
+
204
+ 2
205
+
206
+
207
+
208
+
209
+
210
+ {\displaystyle s={\tfrac {a+b+c}{2}}}
211
+
212
+ is the semiperimeter, or half of the triangle's perimeter.
213
+
214
+ Three other equivalent ways of writing Heron's formula are
215
+
216
+ The area of a parallelogram embedded in a three-dimensional Euclidean space can be calculated using vectors. Let vectors AB and AC point respectively from A to B and from A to C. The area of parallelogram ABDC is then
217
+
218
+ which is the magnitude of the cross product of vectors AB and AC. The area of triangle ABC is half of this,
219
+
220
+ The area of triangle ABC can also be expressed in terms of dot products as follows:
221
+
222
+ In two-dimensional Euclidean space, expressing vector AB as a free vector in Cartesian space equal to (x1,y1) and AC as (x2,y2), this can be rewritten as:
223
+
224
+ If vertex A is located at the origin (0, 0) of a Cartesian coordinate system and the coordinates of the other two vertices are given by B = (xB, yB) and C = (xC, yC), then the area can be computed as ​1⁄2 times the absolute value of the determinant
225
+
226
+ For three general vertices, the equation is:
227
+
228
+ which can be written as
229
+
230
+ If the points are labeled sequentially in the counterclockwise direction, the above determinant expressions are positive and the absolute value signs can be omitted.[13] The above formula is known as the shoelace formula or the surveyor's formula.
231
+
232
+ If we locate the vertices in the complex plane and denote them in counterclockwise sequence as a = xA + yAi, b = xB + yBi, and c = xC + yCi, and denote their complex conjugates as
233
+
234
+
235
+
236
+
237
+
238
+
239
+ a
240
+ ¯
241
+
242
+
243
+
244
+
245
+
246
+ {\displaystyle {\bar {a}}}
247
+
248
+ ,
249
+
250
+
251
+
252
+
253
+
254
+
255
+ b
256
+ ¯
257
+
258
+
259
+
260
+
261
+
262
+ {\displaystyle {\bar {b}}}
263
+
264
+ , and
265
+
266
+
267
+
268
+
269
+
270
+
271
+ c
272
+ ¯
273
+
274
+
275
+
276
+
277
+
278
+ {\displaystyle {\bar {c}}}
279
+
280
+ , then the formula
281
+
282
+ is equivalent to the shoelace formula.
283
+
284
+ In three dimensions, the area of a general triangle A = (xA, yA, zA), B = (xB, yB, zB) and C = (xC, yC, zC) is the Pythagorean sum of the areas of the respective projections on the three principal planes (i.e. x = 0, y = 0 and z = 0):
285
+
286
+ The area within any closed curve, such as a triangle, is given by the line integral around the curve of the algebraic or signed distance of a point on the curve from an arbitrary oriented straight line L. Points to the right of L as oriented are taken to be at negative distance from L, while the weight for the integral is taken to be the component of arc length parallel to L rather than arc length itself.
287
+
288
+ This method is well suited to computation of the area of an arbitrary polygon. Taking L to be the x-axis, the line integral between consecutive vertices (xi,yi) and (xi+1,yi+1) is given by the base times the mean height, namely (xi+1 − xi)(yi + yi+1)/2. The sign of the area is an overall indicator of the direction of traversal, with negative area indicating counterclockwise traversal. The area of a triangle then falls out as the case of a polygon with three sides.
289
+
290
+ While the line integral method has in common with other coordinate-based methods the arbitrary choice of a coordinate system, unlike the others it makes no arbitrary choice of vertex of the triangle as origin or of side as base. Furthermore, the choice of coordinate system defined by L commits to only two degrees of freedom rather than the usual three, since the weight is a local distance (e.g. xi+1 − xi in the above) whence the method does not require choosing an axis normal to L.
291
+
292
+ When working in polar coordinates it is not necessary to convert to Cartesian coordinates to use line integration, since the line integral between consecutive vertices (ri,θi) and (ri+1,θi+1) of a polygon is given directly by riri+1sin(θi+1 − θi)/2. This is valid for all values of θ, with some decrease in numerical accuracy when |θ| is many orders of magnitude greater than π. With this formulation negative area indicates clockwise traversal, which should be kept in mind when mixing polar and cartesian coordinates. Just as the choice of y-axis (x = 0) is immaterial for line integration in cartesian coordinates, so is the choice of zero heading (θ = 0) immaterial here.
293
+
294
+ Three formulas have the same structure as Heron's formula but are expressed in terms of different variables. First, denoting the medians from sides a, b, and c respectively as ma, mb, and mc and their semi-sum (ma + mb + mc)/2 as σ, we have[14]
295
+
296
+ Next, denoting the altitudes from sides a, b, and c respectively as ha, hb, and hc, and denoting the semi-sum of the reciprocals of the altitudes as
297
+
298
+
299
+
300
+ H
301
+ =
302
+ (
303
+
304
+ h
305
+
306
+ a
307
+
308
+
309
+ −
310
+ 1
311
+
312
+
313
+ +
314
+
315
+ h
316
+
317
+ b
318
+
319
+
320
+ −
321
+ 1
322
+
323
+
324
+ +
325
+
326
+ h
327
+
328
+ c
329
+
330
+
331
+ −
332
+ 1
333
+
334
+
335
+ )
336
+
337
+ /
338
+
339
+ 2
340
+
341
+
342
+ {\displaystyle H=(h_{a}^{-1}+h_{b}^{-1}+h_{c}^{-1})/2}
343
+
344
+ we have[15]
345
+
346
+ And denoting the semi-sum of the angles' sines as S = [(sin α) + (sin β) + (sin γ)]/2, we have[16]
347
+
348
+ where D is the diameter of the circumcircle:
349
+
350
+
351
+
352
+ D
353
+ =
354
+
355
+
356
+
357
+ a
358
+
359
+ sin
360
+ ⁡
361
+ α
362
+
363
+
364
+
365
+
366
+ =
367
+
368
+
369
+
370
+ b
371
+
372
+ sin
373
+ ⁡
374
+ β
375
+
376
+
377
+
378
+
379
+ =
380
+
381
+
382
+
383
+ c
384
+
385
+ sin
386
+ ⁡
387
+ γ
388
+
389
+
390
+
391
+
392
+ .
393
+
394
+
395
+ {\displaystyle D={\tfrac {a}{\sin \alpha }}={\tfrac {b}{\sin \beta }}={\tfrac {c}{\sin \gamma }}.}
396
+
397
+ See Pick's theorem for a technique for finding the area of any arbitrary lattice polygon (one drawn on a grid with vertically and horizontally adjacent lattice points at equal distances, and with vertices on lattice points).
398
+
399
+ The theorem states:
400
+
401
+ where
402
+
403
+
404
+
405
+ I
406
+
407
+
408
+ {\displaystyle I}
409
+
410
+ is the number of internal lattice points and B is the number of lattice points lying on the border of the polygon.
411
+
412
+ Numerous other area formulas exist, such as
413
+
414
+ where r is the inradius, and s is the semiperimeter (in fact, this formula holds for all tangential polygons), and[17]:Lemma 2
415
+
416
+ where
417
+
418
+
419
+
420
+
421
+ r
422
+
423
+ a
424
+
425
+
426
+ ,
427
+
428
+
429
+ r
430
+
431
+ b
432
+
433
+
434
+ ,
435
+
436
+
437
+ r
438
+
439
+ c
440
+
441
+
442
+
443
+
444
+ {\displaystyle r_{a},\,r_{b},\,r_{c}}
445
+
446
+ are the radii of the excircles tangent to sides a, b, c respectively.
447
+
448
+ We also have
449
+
450
+ and[18]
451
+
452
+ for circumdiameter D; and[19]
453
+
454
+ for angle α ≠ 90°.
455
+
456
+ The area can also be expressed as[20]
457
+
458
+ In 1885, Baker[21] gave a collection of over a hundred distinct area formulas for the triangle. These include:
459
+
460
+ for circumradius (radius of the circumcircle) R, and
461
+
462
+ The area T of any triangle with perimeter p satisfies
463
+
464
+ with equality holding if and only if the triangle is equilateral.[22][23]:657
465
+
466
+ Other upper bounds on the area T are given by[24]:p.290
467
+
468
+ and
469
+
470
+ both again holding if and only if the triangle is equilateral.
471
+
472
+ There are infinitely many lines that bisect the area of a triangle.[25] Three of them are the medians, which are the only area bisectors that go through the centroid. Three other area bisectors are parallel to the triangle's sides.
473
+
474
+ Any line through a triangle that splits both the triangle's area and its perimeter in half goes through the triangle's incenter. There can be one, two, or three of these for any given triangle.
475
+
476
+ The formulas in this section are true for all Euclidean triangles.
477
+
478
+ The medians and the sides are related by[26]:p.70
479
+
480
+ and
481
+
482
+ and equivalently for mb and mc.
483
+
484
+ For angle A opposite side a, the length of the internal angle bisector is given by[27]
485
+
486
+ for semiperimeter s, where the bisector length is measured from the vertex to where it meets the opposite side.
487
+
488
+ The interior perpendicular bisectors are given by
489
+
490
+ where the sides are
491
+
492
+
493
+
494
+ a
495
+ ≥
496
+ b
497
+ ≥
498
+ c
499
+
500
+
501
+ {\displaystyle a\geq b\geq c}
502
+
503
+ and the area is
504
+
505
+
506
+
507
+ T
508
+ .
509
+
510
+
511
+ {\displaystyle T.}
512
+
513
+ [28]:Thm 2
514
+
515
+ The altitude from, for example, the side of length a is
516
+
517
+ The following formulas involve the circumradius R and the inradius r:
518
+
519
+ where ha etc. are the altitudes to the subscripted sides;[26]:p.79
520
+
521
+ and
522
+
523
+ The product of two sides of a triangle equals the altitude to the third side times the diameter D of the circumcircle:[26]:p.64
524
+
525
+ Suppose two adjacent but non-overlapping triangles share the same side of length f and share the same circumcircle, so that the side of length f is a chord of the circumcircle and the triangles have side lengths (a, b, f) and (c, d, f), with the two triangles together forming a cyclic quadrilateral with side lengths in sequence (a, b, c, d). Then[29]:84
526
+
527
+ Let G be the centroid of a triangle with vertices A, B, and C, and let P be any interior point. Then the distances between the points are related by[29]:174
528
+
529
+ The sum of the squares of the triangle's sides equals three times the sum of the squared distances of the centroid from the vertices:
530
+
531
+ Let qa, qb, and qc be the distances from the centroid to the sides of lengths a, b, and c. Then[29]:173
532
+
533
+ and
534
+
535
+ for area T.
536
+
537
+ Carnot's theorem states that the sum of the distances from the circumcenter to the three sides equals the sum of the circumradius and the inradius.[26]:p.83 Here a segment's length is considered to be negative if and only if the segment lies entirely outside the triangle. This method is especially useful for deducing the properties of more abstract forms of triangles, such as the ones induced by Lie algebras, that otherwise have the same properties as usual triangles.
538
+
539
+ Euler's theorem states that the distance d between the circumcenter and the incenter is given by[26]:p.85
540
+
541
+ or equivalently
542
+
543
+ where R is the circumradius and r is the inradius. Thus for all triangles R ≥ 2r, with equality holding for equilateral triangles.
544
+
545
+ If we denote that the orthocenter divides one altitude into segments of lengths u and v, another altitude into segment lengths w and x, and the third altitude into segment lengths y and z, then uv = wx = yz.[26]:p.94
546
+
547
+ The distance from a side to the circumcenter equals half the distance from the opposite vertex to the orthocenter.[26]:p.99
548
+
549
+ The sum of the squares of the distances from the vertices to the orthocenter H plus the sum of the squares of the sides equals twelve times the square of the circumradius:[26]:p.102
550
+
551
+ In addition to the law of sines, the law of cosines, the law of tangents, and the trigonometric existence conditions given earlier, for any triangle
552
+
553
+ Morley's trisector theorem states that in any triangle, the three points of intersection of the adjacent angle trisectors form an equilateral triangle, called the Morley triangle.
554
+
555
+ As discussed above, every triangle has a unique inscribed circle (incircle) that is interior to the triangle and tangent to all three sides.
556
+
557
+ Every triangle has a unique Steiner inellipse which is interior to the triangle and tangent at the midpoints of the sides. Marden's theorem shows how to find the foci of this ellipse.[31] This ellipse has the greatest area of any ellipse tangent to all three sides of the triangle.
558
+
559
+ The Mandart inellipse of a triangle is the ellipse inscribed within the triangle tangent to its sides at the contact points of its excircles.
560
+
561
+ For any ellipse inscribed in a triangle ABC, let the foci be P and Q. Then[32]
562
+
563
+ Every convex polygon with area T can be inscribed in a triangle of area at most equal to 2T. Equality holds (exclusively) for a parallelogram.[33]
564
+
565
+ The Lemoine hexagon is a cyclic hexagon with vertices given by the six intersections of the sides of a triangle with the three lines that are parallel to the sides and that pass through its symmedian point. In either its simple form or its self-intersecting form, the Lemoine hexagon is interior to the triangle with two vertices on each side of the triangle.
566
+
567
+ Every acute triangle has three inscribed squares (squares in its interior such that all four of a square's vertices lie on a side of the triangle, so two of them lie on the same side and hence one side of the square coincides with part of a side of the triangle). In a right triangle two of the squares coincide and have a vertex at the triangle's right angle, so a right triangle has only two distinct inscribed squares. An obtuse triangle has only one inscribed square, with a side coinciding with part of the triangle's longest side. Within a given triangle, a longer common side is associated with a smaller inscribed square. If an inscribed square has side of length qa and the triangle has a side of length a, part of which side coincides with a side of the square, then qa, a, the altitude ha from the side a, and the triangle's area T are related according to[34][35]
568
+
569
+ The largest possible ratio of the area of the inscribed square to the area of the triangle is 1/2, which occurs when a2 = 2T, q = a/2, and the altitude of the triangle from the base of length a is equal to a. The smallest possible ratio of the side of one inscribed square to the side of another in the same non-obtuse triangle is
570
+
571
+
572
+
573
+ 2
574
+
575
+
576
+ 2
577
+
578
+
579
+
580
+ /
581
+
582
+ 3
583
+ =
584
+ 0.94....
585
+
586
+
587
+ {\displaystyle 2{\sqrt {2}}/3=0.94....}
588
+
589
+ [35] Both of these extreme cases occur for the isosceles right triangle.
590
+
591
+ From an interior point in a reference triangle, the nearest points on the three sides serve as the vertices of the pedal triangle of that point. If the interior point is the circumcenter of the reference triangle, the vertices of the pedal triangle are the midpoints of the reference triangle's sides, and so the pedal triangle is called the midpoint triangle or medial triangle. The midpoint triangle subdivides the reference triangle into four congruent triangles which are similar to the reference triangle.
592
+
593
+ The Gergonne triangle or intouch triangle of a reference triangle has its vertices at the three points of tangency of the reference triangle's sides with its incircle. The extouch triangle of a reference triangle has its vertices at the points of tangency of the reference triangle's excircles with its sides (not extended).
594
+
595
+ The tangential triangle of a reference triangle (other than a right triangle) is the triangle whose sides are on the tangent lines to the reference triangle's circumcircle at its vertices.
596
+
597
+ As mentioned above, every triangle has a unique circumcircle, a circle passing through all three vertices, whose center is the intersection of the perpendicular bisectors of the triangle's sides.
598
+
599
+ Further, every triangle has a unique Steiner circumellipse, which passes through the triangle's vertices and has its center at the triangle's centroid. Of all ellipses going through the triangle's vertices, it has the smallest area.
600
+
601
+ The Kiepert hyperbola is the unique conic which passes through the triangle's three vertices, its centroid, and its circumcenter.
602
+
603
+ Of all triangles contained in a given convex polygon, there exists a triangle with maximal area whose vertices are all vertices of the given polygon.[36]
604
+
605
+ One way to identify locations of points in (or outside) a triangle is to place the triangle in an arbitrary location and orientation in the Cartesian plane, and to use Cartesian coordinates. While convenient for many purposes, this approach has the disadvantage of all points' coordinate values being dependent on the arbitrary placement in the plane.
606
+
607
+ Two systems avoid that feature, so that the coordinates of a point are not affected by moving the triangle, rotating it, or reflecting it as in a mirror, any of which give a congruent triangle, or even by rescaling it to give a similar triangle:
608
+
609
+ A non-planar triangle is a triangle which is not contained in a (flat) plane. Some examples of non-planar triangles in non-Euclidean geometries are spherical triangles in spherical geometry and hyperbolic triangles in hyperbolic geometry.
610
+
611
+ While the measures of the internal angles in planar triangles always sum to 180°, a hyperbolic triangle has measures of angles that sum to less than 180°, and a spherical triangle has measures of angles that sum to more than 180°. A hyperbolic triangle can be obtained by drawing on a negatively curved surface, such as a saddle surface, and a spherical triangle can be obtained by drawing on a positively curved surface such as a sphere. Thus, if one draws a giant triangle on the surface of the Earth, one will find that the sum of the measures of its angles is greater than 180°; in fact it will be between 180° and 540°.[37] In particular it is possible to draw a triangle on a sphere such that the measure of each of its internal angles is equal to 90°, adding up to a total of 270°.
612
+
613
+ Specifically, on a sphere the sum of the angles of a triangle is
614
+
615
+ where f is the fraction of the sphere's area which is enclosed by the triangle. For example, suppose that we draw a triangle on the Earth's surface with vertices at the North Pole, at a point on the equator at 0° longitude, and a point on the equator at 90° West longitude. The great circle line between the latter two points is the equator, and the great circle line between either of those points and the North Pole is a line of longitude; so there are right angles at the two points on the equator. Moreover, the angle at the North Pole is also 90° because the other two vertices differ by 90° of longitude. So the sum of the angles in this triangle is 90° + 90° + 90° = 270°. The triangle encloses 1/4 of the northern hemisphere (90°/360° as viewed from the North Pole) and therefore 1/8 of the Earth's surface, so in the formula f = 1/8; thus the formula correctly gives the sum of the triangle's angles as 270°.
616
+
617
+ From the above angle sum formula we can also see that the Earth's surface is locally flat: If we draw an arbitrarily small triangle in the neighborhood of one point on the Earth's surface, the fraction f of the Earth's surface which is enclosed by the triangle will be arbitrarily close to zero. In this case the angle sum formula simplifies to 180°, which we know is what Euclidean geometry tells us for triangles on a flat surface.
618
+
619
+ Rectangles have been the most popular and common geometric form for buildings since the shape is easy to stack and organize; as a standard, it is easy to design furniture and fixtures to fit inside rectangularly shaped buildings. But triangles, while more difficult to use conceptually, provide a great deal of strength. As computer technology helps architects design creative new buildings, triangular shapes are becoming increasingly prevalent as parts of buildings and as the primary shape for some types of skyscrapers as well as building materials. In Tokyo in 1989, architects had wondered whether it was possible to build a 500-story tower to provide affordable office space for this densely packed city, but with the danger to buildings from earthquakes, architects considered that a triangular shape would be necessary if such a building were to be built.[38]
620
+
621
+ In New York City, as Broadway crisscrosses major avenues, the resulting blocks are cut like triangles, and buildings have been built on these shapes; one such building is the triangularly shaped Flatiron Building which real estate people admit has a "warren of awkward spaces that do not easily accommodate modern office furniture" but that has not prevented the structure from becoming a landmark icon.[39] Designers have made houses in Norway using triangular themes.[40] Triangle shapes have appeared in churches[41] as well as public buildings including colleges[42] as well as supports for innovative home designs.[43]
622
+
623
+ Triangles are sturdy; while a rectangle can collapse into a parallelogram from pressure to one of its points, triangles have a natural strength which supports structures against lateral pressures. A triangle will not change shape unless its sides are bent or extended or broken or if its joints break; in essence, each of the three sides supports the other two. A rectangle, in contrast, is more dependent on the strength of its joints in a structural sense. Some innovative designers have proposed making bricks not out of rectangles, but with triangular shapes which can be combined in three dimensions.[44] It is likely that triangles will be used increasingly in new ways as architecture increases in complexity. It is important to remember that triangles are strong in terms of rigidity, but while packed in a tessellating arrangement triangles are not as strong as hexagons under compression (hence the prevalence of hexagonal forms in nature). Tessellated triangles still maintain superior strength for cantilevering however, and this is the basis for one of the strongest man made structures, the tetrahedral truss.
624
+
625
+
626
+
627
+
628
+
en/579.html.txt ADDED
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1
+
2
+
3
+
4
+
5
+
6
+
7
+
8
+
9
+ Basque (/bæsk, bɑːsk/;[3] Basque: Euskara, [eus̺ˈkaɾa]) is a language spoken in the Basque Country, a region that straddles the westernmost Pyrenees in adjacent parts of Northern Spain and Southwestern France. Linguistically, Basque is unrelated to the other languages of Europe and is a language isolate in relation to any other known living language. The Basques are indigenous to, and primarily inhabit, the Basque Country. The Basque language is spoken by 28.4% (751,500) of Basques in all territories. Of these, 93.2% (700,300) are in the Spanish area of the Basque Country and the remaining 6.8% (51,200) are in the French portion.[1]
10
+
11
+ Native speakers live in a contiguous area that includes parts of four Spanish provinces and the three "ancient provinces" in France. Gipuzkoa, most of Biscay, a few municipalities of Álava and the northern area of Navarre formed the core of the remaining Basque-speaking area before measures were introduced in the 1980s to strengthen the language. By contrast, most of Álava, the western part of Biscay and central and southern areas of Navarre are predominantly populated by native speakers of Spanish, either because Basque was replaced by Spanish over the centuries, in some areas (most of Álava and central Navarre) or because it was possibly never spoken there, in other areas (parts of the Enkarterri and southeastern Navarre).
12
+
13
+ In Francoist Spain, Basque language use was affected by the government’s repressive policies. In the Basque Country, "Francoist repression was not only political, but also linguistic and cultural."[4] The regime placed legal restrictions on the use of language, which was suppressed from official discourse, education, and publishing,[5] making it illegal to register new-born babies under Basque names,[6] and even requiring tombstone engravings in Basque to be removed.[7] In some provinces, the public use of the language was suppressed, with people fined for speaking Basque.[8] Public use of Basque was frowned upon by supporters of the regime, often regarded as a sign of anti-Francoism or separatism.[9] Overall, in the 1960s and later, the trend reversed and education and publishing in Basque began to flourish.[10] As a part of this process, a standardised form of the Basque language, called Euskara Batua, was developed by the Euskaltzaindia in the late 1960s.
14
+
15
+ Besides its standardised version, the five historic Basque dialects are Biscayan, Gipuzkoan and Upper Navarrese in Spain and Navarrese–Lapurdian and Souletin in France. They take their names from the historic Basque provinces, but the dialect boundaries are not congruent with province boundaries. Euskara Batua was created so that the Basque language could be used—and easily understood by all Basque speakers—in formal situations (education, mass media, literature) and this is its main use today. In both Spain and France, the use of Basque for education varies from region to region and from school to school.[11]
16
+
17
+ A language isolate, Basque is believed to be one of the few surviving pre-Indo-European languages in Europe and is the only one in Western Europe. The origin of the Basques and of their languages is not conclusively known, though the most accepted current theory is that early forms of Basque developed before the arrival of Indo-European languages in the area, including the Romance languages that geographically surround the Basque-speaking region. Basque has adopted about 40 percent of its vocabulary from the Romance languages,[12] and Basque speakers have in turn lent their own words to Romance speakers.
18
+
19
+ The Basque alphabet uses the Latin script.
20
+
21
+ In Basque, the name of the language is officially Euskara (alongside various dialect forms).
22
+
23
+ In French, the language is normally called basque, though in recent times euskara has become common. Spanish has a greater variety of names for the language. Today, it is most commonly referred to as el vasco, la lengua vasca, or el euskera. Both terms, vasco and basque, are inherited from the Latin ethnonym Vascones, which in turn goes back to the Greek term οὐασκώνους (ouaskōnous), an ethnonym used by Strabo in his Geographica (23 CE, Book III).[13]
24
+
25
+ The Spanish term Vascuence, derived from Latin vasconĭce,[14] has acquired negative connotations over the centuries and is not well-liked amongst Basque speakers generally. Its use is documented at least as far back as the 14th century when a law passed in Huesca in 1349 stated that Item nuyl corridor nonsia usado que faga mercadería ninguna que compre nin venda entre ningunas personas, faulando en algaravia nin en abraych nin en basquenç: et qui lo fara pague por coto XXX sol—essentially penalising the use of Arabic, Hebrew, or Basque in marketplaces with a fine of 30 sols (the equivalent of 30 sheep).[15]
26
+
27
+ Basque is geographically surrounded by Romance languages but is a language isolate unrelated to them, and indeed, to any other language in the world. It is the last remaining descendant of one of the pre-Indo-European languages of Western Europe, the others being extinct outright.[13] Consequently, its prehistory may not be reconstructible by means of the traditional comparative method except by applying it to differences between dialects within the language. Little is known of its origins, but it is likely that an early form of the Basque language was present in Western Europe before the arrival of the Indo-European languages in the area.
28
+
29
+ Authors such as Miguel de Unamuno and Louis Lucien Bonaparte have noted that the words for "knife" (aizto), "axe" (aizkora), and "hoe" (aitzur) derive from the word for "stone" (haitz), and have therefore concluded that the language dates to prehistoric Europe when those tools were made of stone.[16][17] Others find this unlikely: see the aizkora controversy.
30
+
31
+ Latin inscriptions in Gallia Aquitania preserve a number of words with cognates in the reconstructed proto-Basque language, for instance, the personal names Nescato and Cison (neskato and gizon mean 'young girl' and 'man', respectively in modern Basque). This language is generally referred to as Aquitanian and is assumed to have been spoken in the area before the Roman Republic's conquests in the western Pyrenees. Some authors even argue for late Basquisation, that the language moved westward during Late Antiquity after the fall of the Western Roman Empire into the northern part of Hispania into what is now Basque Country.[13]
32
+
33
+ Roman neglect of this area allowed Aquitanian to survive while the Iberian and Tartessian languages became extinct. Through the long contact with Romance languages, Basque adopted a sizeable number of Romance words. Initially the source was Latin, later Gascon (a branch of Occitan) in the northeast, Navarro-Aragonese in the southeast and Spanish in the southwest.
34
+
35
+ Once accepted as a non-Indo-European language, many attempts have been made to link it with more geographically distant languages. Apart from pseudoscientific comparisons, the appearance of long-range linguistics gave rise to several attempts to connect Basque with geographically very distant language families. Almost all hypotheses concerning the origin of Basque are controversial, and the suggested evidence is not generally accepted by mainstream linguists. Some of these hypothetical connections are:
36
+
37
+ The region where Basque is spoken has become smaller over centuries, especially at the northern, southern, and eastern borders. Nothing is known about the limits of this region in ancient times, but on the basis of toponyms and epigraphs, it seems that in the beginning of the Common Era it stretched to the river Garonne in the north (including the southwestern part of present-day France); at least to the Val d'Aran in the east (now a Gascon-speaking part of Catalonia), including lands on both sides of the Pyrenees;[30] the southern and western boundaries are not clear at all.
38
+
39
+ The Reconquista temporarily counteracted this contracting tendency when the Christian lords called on Northern Iberian peoples—Basques, Asturians, and "Franks"—to colonise the new conquests. The Basque language became the main everyday language,[where?] while other languages like Spanish, Gascon, French, or Latin were preferred for the administration and high education.
40
+
41
+ By the 16th century, the Basque-speaking area was reduced basically to the present-day seven provinces of the Basque Country, excluding the southern part of Navarre, the southwestern part of Álava, and the western part of Biscay, and including some parts of Béarn.[31]
42
+
43
+ In 1807, Basque was still spoken in the northern half of Álava—including its capital city Vitoria-Gasteiz[32]—and a vast area in central Navarre, but in these two provinces, Basque experienced a rapid decline that pushed its border northwards. In the French Basque Country, Basque was still spoken in all the territory except in Bayonne and some villages around, and including some bordering towns in Béarn.
44
+
45
+ In the 20th century, however, the rise of Basque nationalism spurred increased interest in the language as a sign of ethnic identity, and with the establishment of autonomous governments in the Southern Basque Country, it has recently made a modest comeback. In the Spanish part, Basque-language schools for children and Basque-teaching centres for adults have brought the language to areas such as western Enkarterri and the Ribera del Ebro in southern Navarre, where it is not known to ever have been widely spoken; and in the French Basque Country, these schools and centres have almost stopped the decline of the language.
46
+
47
+ Historically, Latin or Romance languages have been the official languages in this region. However, Basque was explicitly recognised in some areas. For instance, the fuero or charter of the Basque-colonised Ojacastro (now in La Rioja) allowed the inhabitants to use Basque in legal processes in the 13th and 14th centuries.
48
+
49
+ The Spanish Constitution of 1978 states in Article 3 that the Spanish language is the official language of the nation, but allows autonomous communities to provide a co-official language status for the other languages of Spain.[33] Consequently, the Statute of Autonomy of the Basque Autonomous Community establishes Basque as the co-official language of the autonomous community. The Statute of Navarre establishes Spanish as the official language of Navarre, but grants co-official status to the Basque language in the Basque-speaking areas of northern Navarre. Basque has no official status in the French Basque Country and French citizens are barred from officially using Basque in a French court of law. However, the use of Basque by Spanish nationals in French courts is permitted (with translation), as Basque is officially recognised on the other side of the border.
50
+
51
+ The positions of the various existing governments differ with regard to the promotion of Basque in areas where Basque is commonly spoken. The language has official status in those territories that are within the Basque Autonomous Community, where it is spoken and promoted heavily, but only partially in Navarre. The Ley del Vascuence ("Law of Basque"), seen as contentious by many Basques, but considered fitting Navarra's linguistic and cultural diversity by some of the main political parties of Navarre,[34] divides Navarre into three language areas: Basque-speaking, non-Basque-speaking, and mixed. Support for the language and the linguistic rights of citizens vary, depending on the area. Others consider it unfair, since the rights of Basque speakers differ greatly depending on the place they live.
52
+
53
+ The 2006 sociolinguistic survey of all Basque-speaking territories showed that in 2006, of all people aged 16 and above:[35]
54
+
55
+ Taken together, in 2006, of a total population of 2,589,600 (1,850,500 in the Autonomous Community, 230,200 in the Northern Provinces and 508,900 in Navarre), 665,800 spoke Basque (aged 16 and above). This amounts to 25.7% Basque bilinguals overall, 15.4% passive speakers, and 58.9% non-speakers. Compared to the 1991 figures, this represents an overall increase of 137,000, from 528,500 (from a population of 2,371,100) 15 years previously.[35]
56
+
57
+ The 2011 figures show an increase of some 64,000 speakers compared to the 2006 figures to 714,136, with significant increases in the Autonomous Community, but a slight drop in the Northern Basque Country to 51,100, overall amounting to an increase to 27% of all inhabitants of Basque provinces (2,648,998 in total).[37]
58
+
59
+ Basque is used as a language of commerce both in the Basque Country and in locations around the world where Basques immigrated throughout history.[40]
60
+
61
+ The modern Basque dialects show a high degree of dialectal divergence, sometimes making cross-dialect communication difficult. This is especially true in the case of Biscayan and Souletin, which are regarded as the most divergent Basque dialects.
62
+
63
+ Modern Basque dialectology distinguishes five dialects:[41]
64
+
65
+ These dialects are divided in 11 subdialects, and 24 minor varieties among them.
66
+ According to Koldo Zuazo,[42] the Biscayan dialect or "Western" is the most widespread dialect, with around 300,000 speakers out of a total of around 660,000 speakers. This dialect is divided in two minor subdialects: the Western Biscayan and Eastern Biscayan, plus transitional dialects.
67
+
68
+ Although the influence of the neighbouring Romance languages on the Basque language (especially the lexicon, but also to some degree Basque phonology and grammar) has been much more extensive, it is usually assumed that there has been some feedback from Basque into these languages as well. In particular Gascon and Aragonese, and to a lesser degree Spanish are thought to have received this influence in the past. In the case of Aragonese and Gascon, this would have been through substrate interference following language shift from Aquitanian or Basque to a Romance language, affecting all levels of the language, including place names around the Pyrenees.[43][44][45][46][47]
69
+
70
+ Although a number of words of alleged Basque origin in the Spanish language are circulated (e.g. anchoa 'anchovies', bizarro 'dashing, gallant, spirited', cachorro 'puppy', etc.), most of these have more easily explicable Romance etymologies or not particularly convincing derivations from Basque.[13] Ignoring cultural terms, there is one strong loanword candidate, ezker, long considered the source of the Pyrennean and Iberian Romance words for "left (side)" (izquierdo, esquerdo, esquerre).[13][48] The lack of initial /r/ in Gascon could arguably be due to a Basque influence but this issue is under-researched.[13]
71
+
72
+ The other most commonly claimed substrate influences:
73
+
74
+ The first two features are common, widespread developments in many Romance (and non-Romance) languages.[13][specify] The change of /f/ to /h/ occurred historically only in a limited area (Gascony and Old Castile) that corresponds almost exactly to areas where heavy Basque bilingualism is assumed, and as a result has been widely postulated (and equally strongly disputed). Substrate theories are often difficult to prove (especially in the case of phonetically plausible changes like /f/ to /h/). As a result, although many arguments have been made on both sides, the debate largely comes down to the a priori tendency on the part of particular linguists to accept or reject substrate arguments.
75
+
76
+ Examples of arguments against the substrate theory,[13] and possible responses:
77
+
78
+ Beyond these arguments, a number of nomadic groups of Castile are also said to use or have used Basque words in their jargon, such as the gacería in Segovia, the mingaña, the Galician fala dos arxinas[49] and the Asturian Xíriga.[50]
79
+
80
+ Part of the Romani community in the Basque Country speaks Erromintxela, which is a rare mixed language, with a Kalderash Romani vocabulary and Basque grammar.[51]
81
+
82
+ A number of Basque-based or Basque-influenced pidgins have existed. In the 16th century, Basque sailors used a Basque–Icelandic pidgin in their contacts with Iceland.[52] The Algonquian–Basque pidgin arose from contact between Basque whalers and the Algonquian peoples in the Gulf of Saint Lawrence and Strait of Belle Isle.[53]
83
+
84
+ Basque is an ergative–absolutive language. The subject of an intransitive verb is in the absolutive case (which is unmarked), and the same case is used for the direct object of a transitive verb. The subject of the transitive verb is marked differently, with the ergative case (shown by the suffix -k). This also triggers main and auxiliary verbal agreement.
85
+
86
+ The auxiliary verb, which accompanies most main verbs, agrees not only with the subject, but with any direct object and the indirect object present. Among European languages, this polypersonal agreement is found only in Basque, some languages of the Caucasus, Mordvinic languages, Hungarian, and Maltese (all non-Indo-European). The ergative–absolutive alignment is also rare among European languages—occurring only in some languages of the Caucasus—but not infrequent worldwide.
87
+
88
+ Consider the phrase:
89
+
90
+ Martin-ek
91
+
92
+ Martin-ERG
93
+
94
+ egunkari-ak
95
+
96
+ newspaper-PL
97
+
98
+ erosten
99
+
100
+ buy-GER
101
+
102
+ di-zki-t
103
+
104
+ AUX.(s)he/it/they.OBJ-PL.OBJ-me.IO[(s)he/it_SBJ]
105
+
106
+ Martin-ek egunkari-ak erosten di-zki-t
107
+
108
+ Martin-ERG newspaper-PL buy-GER AUX.(s)he/it/they.OBJ-PL.OBJ-me.IO[(s)he/it_SBJ]
109
+
110
+ "Martin buys the newspapers for me."
111
+
112
+ Martin-ek is the agent (transitive subject), so it is marked with the ergative case ending -k (with an epenthetic -e-). Egunkariak has an -ak ending, which marks plural object (plural absolutive, direct object case). The verb is erosten dizkitcode: eus promoted to code: eu , in which erostencode: eus promoted to code: eu is a kind of gerund ("buying") and the auxiliary dizkitcode: eus promoted to code: eu means "he/she (does) them for me". This dizkitcode: eus promoted to code: eu can be split like this:
113
+
114
+ The phrase "you buy the newspapers for me" is translated as:
115
+
116
+ Zu-ek
117
+
118
+ you-ERG
119
+
120
+ egunkari-ak
121
+
122
+ newspaper-PL
123
+
124
+ erosten
125
+
126
+ buy-GER
127
+
128
+ di-zki-da-zue
129
+
130
+ AUX.(s)he/it/they.OBJ-PL.OBJ-me.IO-you(pl.).SBJ
131
+
132
+ Zu-ek egunkari-ak erosten di-zki-da-zue
133
+
134
+ you-ERG newspaper-PL buy-GER AUX.(s)he/it/they.OBJ-PL.OBJ-me.IO-you(pl.).SBJ
135
+
136
+
137
+
138
+ The auxiliary verb is composed as di-zki-da-zue and means 'you pl. (do) them for me'
139
+
140
+ The pronoun zuek 'you (plural)' has the same form both in the nominative or absolutive case (the subject of an intransitive sentence or direct object of a transitive sentence) and in the ergative case (the subject of a transitive sentence). In spoken Basque, the auxiliary verb is never dropped even if it is redundant, e.g. dizkidazuecode: eus promoted to code: eu in zuek niri egunkariak erosten dizkidazuecode: eus promoted to code: eu 'you (pl.) are buying the newspapers for me'. However, the pronouns are almost always dropped, e.g. zuek in egunkariak erosten dizkidazuecode: eus promoted to code: eu 'you (pl.) are buying the newspapers for me'. The pronouns are used only to show emphasis: egunkariak zuek erosten dizkidazuecode: eus promoted to code: eu 'it is you (pl.) who buys the newspapers for me', or egunkariak niri erosten dizkidazuecode: eus promoted to code: eu 'it is me for whom you buy the newspapers'.
141
+
142
+ Modern Basque dialects allow for the conjugation of about fifteen verbs, called synthetic verbs, some only in literary contexts. These can be put in the present and past tenses in the indicative and subjunctive moods, in three tenses in the conditional and potential moods, and in one tense in the imperative. Each verb that can be taken intransitively has a nor (absolutive) paradigm and possibly a nor-nori (absolutive–dative) paradigm, as in the sentence Aititeri txapela erori zaio ("The hat fell from grandfather['s head]").[54] Each verb that can be taken transitively uses those two paradigms for antipassive-voice contexts in which no agent is mentioned (Basque lacks a passive voice, and displays instead an antipassive voice paradigm), and also has a nor-nork (absolutive–ergative) paradigm and possibly a nor-nori-nork (absolutive–dative–ergative) paradigm. The last would entail the dizkidazue example above. In each paradigm, each constituent noun can take on any of eight persons, five singular and three plural, with the exception of nor-nori-nork in which the absolutive can only be third person singular or plural. The most ubiquitous auxiliary, izan, can be used in any of these paradigms, depending on the nature of the main verb.
143
+
144
+ There are more persons in the singular (5) than in the plural (3) for synthetic (or filamentous) verbs because of the two familiar persons—informal masculine and feminine second person singular. The pronoun hi is used for both of them, but where the masculine form of the verb uses a -k, the feminine uses an -n. This is a property rarely found in Indo-European languages. The entire paradigm of the verb is further augmented by inflecting for "listener" (the allocutive) even if the verb contains no second person constituent. If the situation calls for the familiar masculine, the form is augmented and modified accordingly. Likewise for the familiar feminine.
145
+ (Gizon bat etorri da, "a man has come"; gizon bat etorri duk, "a man has come [you are a male close friend]", gizon bat etorri dun, "a man has come [you are a female close friend]", gizon bat etorri duzu, "a man has come [I talk to you (Sir / Madam)]")[55] This multiplies the number of possible forms by nearly three. Still, the restriction on contexts in which these forms may be used is strong, since all participants in the conversation must be friends of the same sex, and not too far apart in age. Some dialects dispense with the familiar forms entirely. Note, however, that the formal second person singular conjugates in parallel to the other plural forms, perhaps indicating that it was originally the second person plural, later came to be used as a formal singular, and then later still the modern second person plural was formulated as an innovation.
146
+
147
+ All the other verbs in Basque are called periphrastic, behaving much like a participle would in English. These have only three forms in total, called aspects: perfect (various suffixes), habitual[56] (suffix -t[z]en), and future/potential (suffix. -ko/-go). Verbs of Latinate origin in Basque, as well as many other verbs, have a suffix -tu in the perfect, adapted from the Latin perfect passive -tus suffix. The synthetic verbs also have periphrastic forms, for use in perfects and in simple tenses in which they are deponent.
148
+
149
+ Within a verb phrase, the periphrastic verb comes first, followed by the auxiliary.
150
+
151
+ A Basque noun-phrase is inflected in 17 different ways for case, multiplied by four ways for its definiteness and number (indefinite, definite singular, definite plural, and definite close plural: euskaldun [Basque speaker], euskalduna [the Basque speaker, a Basque speaker], euskaldunak [Basque speakers, the Basque speakers], and euskaldunok [we Basque speakers, those Basque speakers]). These first 68 forms are further modified based on other parts of the sentence, which in turn are inflected for the noun again. It has been estimated that, with two levels of recursion, a Basque noun may have 458,683 inflected forms.[57]
152
+
153
+ The proper name "Mikel" (Michael) is declined as follows:
154
+
155
+ Within a noun phrase, modifying adjectives follow the noun. As an example of a Basque noun phrase, etxe zaharrean "in the old house" is morphologically analysed as follows by Agirre et al.[58]
156
+
157
+ Basic syntactic construction is subject–object–verb (unlike Spanish, French or English where a subject–verb–object construction is more common). The order of the phrases within a sentence can be changed with thematic purposes, whereas the order of the words within a phrase is usually rigid. As a matter of fact, Basque phrase order is topic–focus, meaning that in neutral sentences (such as sentences to inform someone of a fact or event) the topic is stated first, then the focus. In such sentences, the verb phrase comes at the end. In brief, the focus directly precedes the verb phrase. This rule is also applied in questions, for instance, What is this? can be translated as Zer da hau? or Hau zer da?, but in both cases the question tag zer immediately precedes the verb da. This rule is so important in Basque that, even in grammatical descriptions of Basque in other languages, the Basque word galdegai (focus) is used.[clarification needed]
158
+
159
+ In negative sentences, the order changes. Since the negative particle ez must always directly precede the auxiliary, the topic most often comes beforehand, and the rest of the sentence follows. This includes the periphrastic, if there is one: Aitak frantsesa irakasten du, "Father teaches French," in the negative becomes Aitak ez du frantsesa irakasten, in which irakasten ("teaching") is separated from its auxiliary and placed at the end.
160
+
161
+ The Basque language features five vowels: /a/, /e/, /i/, /o/ and /u/ (the same that are found in Spanish, Asturian and Aragonese). In the Zuberoan dialect, extra phonemes are featured:
162
+
163
+ Basque has a distinction between laminal and apical articulation for the alveolar fricatives and affricates. With the laminal alveolar fricative [s̻], the friction occurs across the blade of the tongue, the tongue tip pointing toward the lower teeth. This is the usual /s/ in most European languages. It is written with an orthographic ⟨z⟩. By contrast, the voiceless apicoalveolar fricative [s̺] is written ⟨s⟩; the tip of the tongue points toward the upper teeth and friction occurs at the tip (apex). For example, zu "you" (singular, respectful) is distinguished from su "fire". The affricate counterparts are written ⟨tz⟩ and ⟨ts⟩. So, etzi "the day after tomorrow" is distinguished from etsi "to give up"; atzo "yesterday" is distinguished from atso "old woman".
164
+
165
+ In the westernmost parts of the Basque country, only the apical ⟨s⟩ and the alveolar affricate ⟨tz⟩ are used.
166
+
167
+ Basque also features postalveolar sibilants (/ʃ/, written ⟨x⟩, and /tʃ/, written ⟨tx⟩), sounding like English sh and ch.
168
+
169
+ There are two palatal stops, voiced and unvoiced, as well as a palatal nasal and a palatal lateral (the palatal stops are not present in all dialects). These and the postalveolar sounds are typical of diminutives, which are used frequently in child language and motherese (mainly to show affection rather than size). For example, tanta "drop" vs. ttantta /canca/ "droplet". A few common words, such as txakur /tʃakur/ "dog", use palatal sounds even though in current usage they have lost the diminutive sense, the corresponding non-palatal forms now acquiring an augmentative or pejorative sense: zakur—"big dog". Many Basque dialects exhibit a derived palatalisation effect, in which coronal onset consonants change into the palatal counterpart after the high front vowel /i/. For example, the /n/ in egin "to act" becomes palatal in southern and western dialects when a suffix beginning with a vowel is added: /eɡina/ = [eɡiɲa] "the action", /eɡines̻/ = [eɡiɲes̻] "doing".
170
+
171
+ The letter ⟨j⟩ has a variety of realisations according to the regional dialect: [j, dʒ, x, ʃ, ɟ, ʝ], as pronounced from west to east in south Bizkaia and coastal Lapurdi, central Bizkaia, east Bizkaia and Gipuzkoa, south Navarre, inland Lapurdi and Low Navarre, and Zuberoa, respectively.[59]
172
+
173
+ The letter ⟨h⟩ is silent in the Southern dialects, but pronounced (although vanishing) in the Northern ones. Unified Basque spells it except when it is predictable, in a position following a consonant.[clarification needed][60]
174
+
175
+ Unless they are recent loanwords (e.g. Ruanda (Rwanda), radar ... ), words may not have initial ⟨r⟩. In older loans, initial r- took a prosthetic e-, resulting in err- (Erroma "Rome", Errusia "Russia"), more rarely irr- (for example irratia "radio", irrisa "rice").
176
+
177
+ Basque features great dialectal variation in accentuation, from a weak pitch accent in the western dialects to a marked stress in central and eastern dialects, with varying patterns of stress placement.[61] Stress is in general not distinctive (and for historical comparisons not very useful); there are, however, a few instances where stress is phonemic, serving to distinguish between a few pairs of stress-marked words and between some grammatical forms (mainly plurals from other forms), e.g. basóà ("the forest", absolutive case) vs. básoà ("the glass", absolutive case; an adoption from Spanish vaso); basóàk ("the forest", ergative case) vs. básoàk ("the glass", ergative case) vs. básoak ("the forests" or "the glasses", absolutive case).
178
+
179
+ Given its great deal of variation among dialects, stress is not marked in the standard orthography and Euskaltzaindia (the Academy of the Basque Language) provides only general recommendations for a standard placement of stress, basically to place a high-pitched weak stress (weaker than that of Spanish, let alone that of English) on the second syllable of a syntagma, and a low-pitched even-weaker stress on its last syllable, except in plural forms where stress is moved to the first syllable.
180
+
181
+ This scheme provides Basque with a distinct musicality[citation needed] that differentiates its sound from the prosodical patterns of Spanish (which tends to stress the second-to-last syllable). Some Euskaldun berriak ("new Basque-speakers", i.e. second-language Basque-speakers) with Spanish as their first language tend to carry the prosodical patterns of Spanish into their pronunciation of Basque, e.g. pronouncing nire ama ("my mum") as nire áma (– – ´ –), instead of as niré amà (– ´ – `).
182
+
183
+ The combining forms of nominals in final /-u/ vary across the regions of the Basque Country. The /u/ can stay unchanged, be lowered to an /a/, or it can be lost. Loss is most common in the east, while lowering is most common in the west. For instance, buru, "head", has the combining forms buru- and bur-, as in buruko, "cap", and burko, "pillow", whereas katu, "cat", has the combining form kata-, as in katakume, "kitten". Michelena suggests that the lowering to /a/ is generalised from cases of Romance borrowings in Basque that retained Romance stem alternations, such as kantu, "song" with combining form kanta-, borrowed from Romance canto, canta-.[62]
184
+
185
+ Through contact with neighbouring peoples, Basque has adopted many words from Latin, Spanish, and Gascon, among other languages. There are a considerable number of Latin loans (sometimes obscured by being subject to Basque phonology and grammar for centuries), for example: lore ("flower", from florem), errota ("mill", from rotam, "[mill] wheel"), gela ("room", from cellam), gauza ("thing", from causa).
186
+
187
+ Basque is written using the Latin script including ñ and sometimes ç and ü. Basque does not use Cc, Qq, Vv, Ww, Yy for words that have some tradition in this language; nevertheless, the Basque alphabet (established by Euskaltzaindia) does include them for loanwords:[63]
188
+
189
+ The phonetically meaningful digraphs dd, ll, rr, ts, tt, tx, tz are treated as pairs of letters.
190
+
191
+ All letters and digraphs represent unique phonemes. The main exception is when l and n are preceded by i, that in most dialects palatalises their sound into /ʎ/ and /ɲ/, even if these are not written. Hence, Ikurriña can also be written Ikurrina without changing the sound, whereas the proper name Ainhoa requires the mute h to break the palatalisation of the n.
192
+
193
+ H is mute in most regions, but it is pronounced in many places in the northeast, the main reason for its existence in the Basque alphabet.
194
+ Its acceptance was a matter of contention during the standardisation process because the speakers of the most extended dialects had to learn where to place these h's, silent for them.
195
+
196
+ In Sabino Arana's (1865–1903) alphabet,[64] digraphs ⟨ll⟩ and ⟨rr⟩ were replaced with ĺ and ŕ, respectively.
197
+
198
+ A typically Basque style of lettering is sometimes used for inscriptions.
199
+ It derives from the work of stone and wood carvers and is characterised by thick serifs.
200
+
201
+ Basque millers traditionally employed a separate number system of unknown origin.[65] In this system the symbols are arranged either along a vertical line or horizontally. On the vertical line the single digits and fractions are usually off to one side, usually at the top. When used horizontally, the smallest units are usually on the right and the largest on the left.
202
+
203
+ The system is, as is the Basque system of counting in general, vigesimal (base 20). Although the system is in theory capable of indicating numbers above 100, most recorded examples do not go above 100 in general. Fractions are relatively common, especially ​1⁄2.
204
+
205
+ The exact systems used vary from area to area but generally follow the same principle with 5 usually being a diagonal line or a curve off the vertical line (a V shape is used when writing a 5 horizontally). Units of ten are usually a horizontal line through the vertical. The twenties are based on a circle with intersecting lines. This system is no longer in general use but is occasionally employed for decorative purposes.
206
+
207
+ Esklabu erremintaria
208
+ Sartaldeko oihanetan gatibaturik
209
+ Erromara ekarri zinduten, esklabua,
210
+ erremintari ofizioa eman zizuten
211
+ eta kateak egiten dituzu.
212
+ Labetik ateratzen duzun burdin goria
213
+ nahieran molda zenezake,
214
+ ezpatak egin ditzakezu
215
+ zure herritarrek kateak hauts ditzaten,
216
+ baina zuk, esklabu horrek,
217
+ kateak egiten dituzu, kate gehiago.
218
+
219
+ IPA pronunciation
220
+ [s̺artaldeko ojanetan ɡatibatuɾik
221
+ eromaɾa ekari s̻induten es̺klabua
222
+ eremintaɾi ofis̻ioa eman s̻is̻uten
223
+ eta kateak eɡiten ditus̻u
224
+ labetik ateɾats̻en dus̻un burdiɲ ɡoɾia
225
+ najeɾan molda s̻enes̻ake
226
+ es̻patak eɡin dits̻akes̻u
227
+ s̻uɾe eritarek kateak auts̺ dits̻aten
228
+ baɲa s̻uk es̺klabu orek
229
+ kateak eɡiten ditus̻u kate ɡejaɡo]
230
+
231
+ The blacksmith slave
232
+ Captive in the rainforests of the West
233
+ they brought you to Rome, slave,
234
+ they gave you the blacksmith work
235
+ and you make chains.
236
+ The incandescent iron you take out of the oven
237
+ can be adapted as you wish,
238
+ you could make swords
239
+ so your people could break the chains,
240
+ but you, o, slave,
241
+ you make chains, more chains.
en/5790.html.txt ADDED
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1
+
2
+
3
+ A triangle is a polygon with three edges and three vertices. It is one of the basic shapes in geometry. A triangle with vertices A, B, and C is denoted
4
+
5
+
6
+
7
+ △
8
+ A
9
+ B
10
+ C
11
+
12
+
13
+ {\displaystyle \triangle ABC}
14
+
15
+ .
16
+
17
+ In Euclidean geometry any three points, when non-collinear, determine a unique triangle and simultaneously, a unique plane (i.e. a two-dimensional Euclidean space). In other words, there is only one plane that contains that triangle, and every triangle is contained in some plane. If the entire geometry is only the Euclidean plane, there is only one plane and all triangles are contained in it; however, in higher-dimensional Euclidean spaces, this is no longer true. This article is about triangles in Euclidean geometry, and in particular, the Euclidean plane, except where otherwise noted.
18
+
19
+ Triangles can be classified according to the lengths of their sides:
20
+
21
+ Hatch marks, also called tick marks, are used in diagrams of triangles and other geometric figures to identify sides of equal lengths. A side can be marked with a pattern of "ticks", short line segments in the form of tally marks; two sides have equal lengths if they are both marked with the same pattern. In a triangle, the pattern is usually no more than 3 ticks. An equilateral triangle has the same pattern on all 3 sides, an isosceles triangle has the same pattern on just 2 sides, and a scalene triangle has different patterns on all sides since no sides are equal. Similarly, patterns of 1, 2, or 3 concentric arcs inside the angles are used to indicate equal angles. An equilateral triangle has the same pattern on all 3 angles, an isosceles triangle has the same pattern on just 2 angles, and a scalene triangle has different patterns on all angles since no angles are equal.
22
+
23
+ Triangles can also be classified according to their internal angles, measured here in degrees.
24
+
25
+ A triangle that has two angles with the same measure also has two sides with the same length, and therefore it is an isosceles triangle. It follows that in a triangle where all angles have the same measure, all three sides have the same length, and such a triangle is therefore equilateral.
26
+
27
+ Triangles are assumed to be two-dimensional plane figures, unless the context provides otherwise (see Non-planar triangles, below). In rigorous treatments, a triangle is therefore called a 2-simplex (see also Polytope). Elementary facts about triangles were presented by Euclid in books 1–4 of his Elements, around 300 BC.
28
+
29
+ The sum of the measures of the interior angles of a triangle in Euclidean space is always 180 degrees.[5] This fact is equivalent to Euclid's parallel postulate. This allows determination of the measure of the third angle of any triangle given the measure of two angles. An exterior angle of a triangle is an angle that is a linear pair (and hence supplementary) to an interior angle. The measure of an exterior angle of a triangle is equal to the sum of the measures of the two interior angles that are not adjacent to it; this is the exterior angle theorem. The sum of the measures of the three exterior angles (one for each vertex) of any triangle is 360 degrees.[note 2]
30
+
31
+ Two triangles are said to be similar if every angle of one triangle has the same measure as the corresponding angle in the other triangle. The corresponding sides of similar triangles have lengths that are in the same proportion, and this property is also sufficient to establish similarity.
32
+
33
+ Some basic theorems about similar triangles are:
34
+
35
+ Two triangles that are congruent have exactly the same size and shape:[note 4] all pairs of corresponding interior angles are equal in measure, and all pairs of corresponding sides have the same length. (This is a total of six equalities, but three are often sufficient to prove congruence.)
36
+
37
+ Some individually necessary and sufficient conditions for a pair of triangles to be congruent are:
38
+
39
+ Some individually sufficient conditions are:
40
+
41
+ An important condition is:
42
+
43
+ Using right triangles and the concept of similarity, the trigonometric functions sine and cosine can be defined. These are functions of an angle which are investigated in trigonometry.
44
+
45
+ A central theorem is the Pythagorean theorem, which states in any right triangle, the square of the length of the hypotenuse equals the sum of the squares of the lengths of the two other sides. If the hypotenuse has length c, and the legs have lengths a and b, then the theorem states that
46
+
47
+ The converse is true: if the lengths of the sides of a triangle satisfy the above equation, then the triangle has a right angle opposite side c.
48
+
49
+ Some other facts about right triangles:
50
+
51
+ For all triangles, angles and sides are related by the law of cosines and law of sines (also called the cosine rule and sine rule).
52
+
53
+ The triangle inequality states that the sum of the lengths of any two sides of a triangle must be greater than or equal to the length of the third side. That sum can equal the length of the third side only in the case of a degenerate triangle, one with collinear vertices. It is not possible for that sum to be less than the length of the third side. A triangle with three given positive side lengths exists if and only if those side lengths satisfy the triangle inequality.
54
+
55
+ Three given angles form a non-degenerate triangle (and indeed an infinitude of them) if and only if both of these conditions hold: (a) each of the angles is positive, and (b) the angles sum to 180°. If degenerate triangles are permitted, angles of 0° are permitted.
56
+
57
+ Three positive angles α, β, and γ, each of them less than 180°, are the angles of a triangle if and only if any one of the following conditions holds:
58
+
59
+ the last equality applying only if none of the angles is 90° (so the tangent function's value is always finite).
60
+
61
+ There are thousands of different constructions that find a special point associated with (and often inside) a triangle, satisfying some unique property: see the article Encyclopedia of Triangle Centers for a catalogue of them. Often they are constructed by finding three lines associated in a symmetrical way with the three sides (or vertices) and then proving that the three lines meet in a single point: an important tool for proving the existence of these is Ceva's theorem, which gives a criterion for determining when three such lines are concurrent. Similarly, lines associated with a triangle are often constructed by proving that three symmetrically constructed points are collinear: here Menelaus' theorem gives a useful general criterion. In this section just a few of the most commonly encountered constructions are explained.
62
+
63
+ A perpendicular bisector of a side of a triangle is a straight line passing through the midpoint of the side and being perpendicular to it, i.e. forming a right angle with it. The three perpendicular bisectors meet in a single point, the triangle's circumcenter, usually denoted by O; this point is the center of the circumcircle, the circle passing through all three vertices. The diameter of this circle, called the circumdiameter, can be found from the law of sines stated above. The circumcircle's radius is called the circumradius.
64
+
65
+ Thales' theorem implies that if the circumcenter is located on a side of the triangle, then the opposite angle is a right one. If the circumcenter is located inside the triangle, then the triangle is acute; if the circumcenter is located outside the triangle, then the triangle is obtuse.
66
+
67
+ An altitude of a triangle is a straight line through a vertex and perpendicular to (i.e. forming a right angle with) the opposite side. This opposite side is called the base of the altitude, and the point where the altitude intersects the base (or its extension) is called the foot of the altitude. The length of the altitude is the distance between the base and the vertex. The three altitudes intersect in a single point, called the orthocenter of the triangle, usually denoted by H. The orthocenter lies inside the triangle if and only if the triangle is acute.
68
+
69
+ An angle bisector of a triangle is a straight line through a vertex which cuts the corresponding angle in half. The three angle bisectors intersect in a single point, the incenter, usually denoted by I, the center of the triangle's incircle. The incircle is the circle which lies inside the triangle and touches all three sides. Its radius is called the inradius. There are three other important circles, the excircles; they lie outside the triangle and touch one side as well as the extensions of the other two. The centers of the in- and excircles form an orthocentric system.
70
+
71
+ A median of a triangle is a straight line through a vertex and the midpoint of the opposite side, and divides the triangle into two equal areas. The three medians intersect in a single point, the triangle's centroid or geometric barycenter, usually denoted by G. The centroid of a rigid triangular object (cut out of a thin sheet of uniform density) is also its center of mass: the object can be balanced on its centroid in a uniform gravitational field. The centroid cuts every median in the ratio 2:1, i.e. the distance between a vertex and the centroid is twice the distance between the centroid and the midpoint of the opposite side.
72
+
73
+ The midpoints of the three sides and the feet of the three altitudes all lie on a single circle, the triangle's nine-point circle. The remaining three points for which it is named are the midpoints of the portion of altitude between the vertices and the orthocenter. The radius of the nine-point circle is half that of the circumcircle. It touches the incircle (at the Feuerbach point) and the three excircles.
74
+
75
+ The orthocenter (blue point), center of the nine-point circle (red), centroid (orange), and circumcenter (green) all lie on a single line, known as Euler's line (red line). The center of the nine-point circle lies at the midpoint between the orthocenter and the circumcenter, and the distance between the centroid and the circumcenter is half that between the centroid and the orthocenter.
76
+
77
+ The center of the incircle is not in general located on Euler's line.
78
+
79
+ If one reflects a median in the angle bisector that passes through the same vertex, one obtains a symmedian. The three symmedians intersect in a single point, the symmedian point of the triangle.
80
+
81
+ There are various standard methods for calculating the length of a side or the measure of an angle. Certain methods are suited to calculating values in a right-angled triangle; more complex methods may be required in other situations.
82
+
83
+ In right triangles, the trigonometric ratios of sine, cosine and tangent can be used to find unknown angles and the lengths of unknown sides. The sides of the triangle are known as follows:
84
+
85
+ The sine of an angle is the ratio of the length of the opposite side to the length of the hypotenuse. In our case
86
+
87
+ This ratio does not depend on the particular right triangle chosen, as long as it contains the angle A, since all those triangles are similar.
88
+
89
+ The cosine of an angle is the ratio of the length of the adjacent side to the length of the hypotenuse. In our case
90
+
91
+ The tangent of an angle is the ratio of the length of the opposite side to the length of the adjacent side. In our case
92
+
93
+ The acronym "SOH-CAH-TOA" is a useful mnemonic for these ratios.
94
+
95
+ The inverse trigonometric functions can be used to calculate the internal angles for a right angled triangle with the length of any two sides.
96
+
97
+ Arcsin can be used to calculate an angle from the length of the opposite side and the length of the hypotenuse.
98
+
99
+ Arccos can be used to calculate an angle from the length of the adjacent side and the length of the hypotenuse.
100
+
101
+ Arctan can be used to calculate an angle from the length of the opposite side and the length of the adjacent side.
102
+
103
+ In introductory geometry and trigonometry courses, the notation sin−1, cos−1, etc., are often used in place of arcsin, arccos, etc. However, the arcsin, arccos, etc., notation is standard in higher mathematics where trigonometric functions are commonly raised to powers, as this avoids confusion between multiplicative inverse and compositional inverse.
104
+
105
+ The law of sines, or sine rule,[8] states that the ratio of the length of a side to the sine of its corresponding opposite angle is constant, that is
106
+
107
+ This ratio is equal to the diameter of the circumscribed circle of the given triangle. Another interpretation of this theorem is that every triangle with angles α, β and γ is similar to a triangle with side lengths equal to sin α, sin β and sin γ. This triangle can be constructed by first constructing a circle of diameter 1, and inscribing in it two of the angles of the triangle. The length of the sides of that triangle will be sin α, sin β and sin γ. The side whose length is sin α is opposite to the angle whose measure is α, etc.
108
+
109
+ The law of cosines, or cosine rule, connects the length of an unknown side of a triangle to the length of the other sides and the angle opposite to the unknown side.[8] As per the law:
110
+
111
+ For a triangle with length of sides a, b, c and angles of α, β, γ respectively, given two known lengths of a triangle a and b, and the angle between the two known sides γ (or the angle opposite to the unknown side c), to calculate the third side c, the following formula can be used:
112
+
113
+ If the lengths of all three sides of any triangle are known the three angles can be calculated:
114
+
115
+ The law of tangents, or tangent rule, can be used to find a side or an angle when two sides and an angle or two angles and a side are known. It states that:[9]
116
+
117
+ "Solution of triangles" is the main trigonometric problem: to find missing characteristics of a triangle (three angles, the lengths of the three sides etc.) when at least three of these characteristics are given. The triangle can be located on a plane or on a sphere. This problem often occurs in various trigonometric applications, such as geodesy, astronomy, construction, navigation etc.
118
+
119
+ Calculating the area T of a triangle is an elementary problem encountered often in many different situations. The best known and simplest formula is:
120
+
121
+ where b is the length of the base of the triangle, and h is the height or altitude of the triangle. The term "base" denotes any side, and "height" denotes the length of a perpendicular from the vertex opposite the base onto the line containing the base. In 499 CE Aryabhata, used this illustrated method in the Aryabhatiya (section 2.6).[10]
122
+
123
+ Although simple, this formula is only useful if the height can be readily found, which is not always the case. For example, the surveyor of a triangular field might find it relatively easy to measure the length of each side, but relatively difficult to construct a 'height'. Various methods may be used in practice, depending on what is known about the triangle. The following is a selection of frequently used formulae for the area of a triangle.[11]
124
+
125
+ The height of a triangle can be found through the application of trigonometry.
126
+
127
+ Knowing SAS: Using the labels in the image on the right, the altitude is h = a sin
128
+
129
+
130
+
131
+ γ
132
+
133
+
134
+ {\displaystyle \gamma }
135
+
136
+ . Substituting this in the formula
137
+
138
+
139
+
140
+ T
141
+ =
142
+
143
+
144
+ 1
145
+ 2
146
+
147
+
148
+ b
149
+ h
150
+
151
+
152
+ {\displaystyle T={\frac {1}{2}}bh}
153
+
154
+ derived above, the area of the triangle can be expressed as:
155
+
156
+ (where α is the interior angle at A, β is the interior angle at B,
157
+
158
+
159
+
160
+ γ
161
+
162
+
163
+ {\displaystyle \gamma }
164
+
165
+ is the interior angle at C and c is the line AB).
166
+
167
+ Furthermore, since sin α = sin (π − α) = sin (β +
168
+
169
+
170
+
171
+ γ
172
+
173
+
174
+ {\displaystyle \gamma }
175
+
176
+ ), and similarly for the other two angles:
177
+
178
+ Knowing AAS:
179
+
180
+ and analogously if the known side is a or c.
181
+
182
+ Knowing ASA:[12]
183
+
184
+ and analogously if the known side is b or c.
185
+
186
+ The shape of the triangle is determined by the lengths of the sides. Therefore, the area can also be derived from the lengths of the sides. By Heron's formula:
187
+
188
+ where
189
+
190
+
191
+
192
+ s
193
+ =
194
+
195
+
196
+
197
+
198
+ a
199
+ +
200
+ b
201
+ +
202
+ c
203
+
204
+ 2
205
+
206
+
207
+
208
+
209
+
210
+ {\displaystyle s={\tfrac {a+b+c}{2}}}
211
+
212
+ is the semiperimeter, or half of the triangle's perimeter.
213
+
214
+ Three other equivalent ways of writing Heron's formula are
215
+
216
+ The area of a parallelogram embedded in a three-dimensional Euclidean space can be calculated using vectors. Let vectors AB and AC point respectively from A to B and from A to C. The area of parallelogram ABDC is then
217
+
218
+ which is the magnitude of the cross product of vectors AB and AC. The area of triangle ABC is half of this,
219
+
220
+ The area of triangle ABC can also be expressed in terms of dot products as follows:
221
+
222
+ In two-dimensional Euclidean space, expressing vector AB as a free vector in Cartesian space equal to (x1,y1) and AC as (x2,y2), this can be rewritten as:
223
+
224
+ If vertex A is located at the origin (0, 0) of a Cartesian coordinate system and the coordinates of the other two vertices are given by B = (xB, yB) and C = (xC, yC), then the area can be computed as ​1⁄2 times the absolute value of the determinant
225
+
226
+ For three general vertices, the equation is:
227
+
228
+ which can be written as
229
+
230
+ If the points are labeled sequentially in the counterclockwise direction, the above determinant expressions are positive and the absolute value signs can be omitted.[13] The above formula is known as the shoelace formula or the surveyor's formula.
231
+
232
+ If we locate the vertices in the complex plane and denote them in counterclockwise sequence as a = xA + yAi, b = xB + yBi, and c = xC + yCi, and denote their complex conjugates as
233
+
234
+
235
+
236
+
237
+
238
+
239
+ a
240
+ ¯
241
+
242
+
243
+
244
+
245
+
246
+ {\displaystyle {\bar {a}}}
247
+
248
+ ,
249
+
250
+
251
+
252
+
253
+
254
+
255
+ b
256
+ ¯
257
+
258
+
259
+
260
+
261
+
262
+ {\displaystyle {\bar {b}}}
263
+
264
+ , and
265
+
266
+
267
+
268
+
269
+
270
+
271
+ c
272
+ ¯
273
+
274
+
275
+
276
+
277
+
278
+ {\displaystyle {\bar {c}}}
279
+
280
+ , then the formula
281
+
282
+ is equivalent to the shoelace formula.
283
+
284
+ In three dimensions, the area of a general triangle A = (xA, yA, zA), B = (xB, yB, zB) and C = (xC, yC, zC) is the Pythagorean sum of the areas of the respective projections on the three principal planes (i.e. x = 0, y = 0 and z = 0):
285
+
286
+ The area within any closed curve, such as a triangle, is given by the line integral around the curve of the algebraic or signed distance of a point on the curve from an arbitrary oriented straight line L. Points to the right of L as oriented are taken to be at negative distance from L, while the weight for the integral is taken to be the component of arc length parallel to L rather than arc length itself.
287
+
288
+ This method is well suited to computation of the area of an arbitrary polygon. Taking L to be the x-axis, the line integral between consecutive vertices (xi,yi) and (xi+1,yi+1) is given by the base times the mean height, namely (xi+1 − xi)(yi + yi+1)/2. The sign of the area is an overall indicator of the direction of traversal, with negative area indicating counterclockwise traversal. The area of a triangle then falls out as the case of a polygon with three sides.
289
+
290
+ While the line integral method has in common with other coordinate-based methods the arbitrary choice of a coordinate system, unlike the others it makes no arbitrary choice of vertex of the triangle as origin or of side as base. Furthermore, the choice of coordinate system defined by L commits to only two degrees of freedom rather than the usual three, since the weight is a local distance (e.g. xi+1 − xi in the above) whence the method does not require choosing an axis normal to L.
291
+
292
+ When working in polar coordinates it is not necessary to convert to Cartesian coordinates to use line integration, since the line integral between consecutive vertices (ri,θi) and (ri+1,θi+1) of a polygon is given directly by riri+1sin(θi+1 − θi)/2. This is valid for all values of θ, with some decrease in numerical accuracy when |θ| is many orders of magnitude greater than π. With this formulation negative area indicates clockwise traversal, which should be kept in mind when mixing polar and cartesian coordinates. Just as the choice of y-axis (x = 0) is immaterial for line integration in cartesian coordinates, so is the choice of zero heading (θ = 0) immaterial here.
293
+
294
+ Three formulas have the same structure as Heron's formula but are expressed in terms of different variables. First, denoting the medians from sides a, b, and c respectively as ma, mb, and mc and their semi-sum (ma + mb + mc)/2 as σ, we have[14]
295
+
296
+ Next, denoting the altitudes from sides a, b, and c respectively as ha, hb, and hc, and denoting the semi-sum of the reciprocals of the altitudes as
297
+
298
+
299
+
300
+ H
301
+ =
302
+ (
303
+
304
+ h
305
+
306
+ a
307
+
308
+
309
+ −
310
+ 1
311
+
312
+
313
+ +
314
+
315
+ h
316
+
317
+ b
318
+
319
+
320
+ −
321
+ 1
322
+
323
+
324
+ +
325
+
326
+ h
327
+
328
+ c
329
+
330
+
331
+ −
332
+ 1
333
+
334
+
335
+ )
336
+
337
+ /
338
+
339
+ 2
340
+
341
+
342
+ {\displaystyle H=(h_{a}^{-1}+h_{b}^{-1}+h_{c}^{-1})/2}
343
+
344
+ we have[15]
345
+
346
+ And denoting the semi-sum of the angles' sines as S = [(sin α) + (sin β) + (sin γ)]/2, we have[16]
347
+
348
+ where D is the diameter of the circumcircle:
349
+
350
+
351
+
352
+ D
353
+ =
354
+
355
+
356
+
357
+ a
358
+
359
+ sin
360
+ ⁡
361
+ α
362
+
363
+
364
+
365
+
366
+ =
367
+
368
+
369
+
370
+ b
371
+
372
+ sin
373
+ ⁡
374
+ β
375
+
376
+
377
+
378
+
379
+ =
380
+
381
+
382
+
383
+ c
384
+
385
+ sin
386
+ ⁡
387
+ γ
388
+
389
+
390
+
391
+
392
+ .
393
+
394
+
395
+ {\displaystyle D={\tfrac {a}{\sin \alpha }}={\tfrac {b}{\sin \beta }}={\tfrac {c}{\sin \gamma }}.}
396
+
397
+ See Pick's theorem for a technique for finding the area of any arbitrary lattice polygon (one drawn on a grid with vertically and horizontally adjacent lattice points at equal distances, and with vertices on lattice points).
398
+
399
+ The theorem states:
400
+
401
+ where
402
+
403
+
404
+
405
+ I
406
+
407
+
408
+ {\displaystyle I}
409
+
410
+ is the number of internal lattice points and B is the number of lattice points lying on the border of the polygon.
411
+
412
+ Numerous other area formulas exist, such as
413
+
414
+ where r is the inradius, and s is the semiperimeter (in fact, this formula holds for all tangential polygons), and[17]:Lemma 2
415
+
416
+ where
417
+
418
+
419
+
420
+
421
+ r
422
+
423
+ a
424
+
425
+
426
+ ,
427
+
428
+
429
+ r
430
+
431
+ b
432
+
433
+
434
+ ,
435
+
436
+
437
+ r
438
+
439
+ c
440
+
441
+
442
+
443
+
444
+ {\displaystyle r_{a},\,r_{b},\,r_{c}}
445
+
446
+ are the radii of the excircles tangent to sides a, b, c respectively.
447
+
448
+ We also have
449
+
450
+ and[18]
451
+
452
+ for circumdiameter D; and[19]
453
+
454
+ for angle α ≠ 90°.
455
+
456
+ The area can also be expressed as[20]
457
+
458
+ In 1885, Baker[21] gave a collection of over a hundred distinct area formulas for the triangle. These include:
459
+
460
+ for circumradius (radius of the circumcircle) R, and
461
+
462
+ The area T of any triangle with perimeter p satisfies
463
+
464
+ with equality holding if and only if the triangle is equilateral.[22][23]:657
465
+
466
+ Other upper bounds on the area T are given by[24]:p.290
467
+
468
+ and
469
+
470
+ both again holding if and only if the triangle is equilateral.
471
+
472
+ There are infinitely many lines that bisect the area of a triangle.[25] Three of them are the medians, which are the only area bisectors that go through the centroid. Three other area bisectors are parallel to the triangle's sides.
473
+
474
+ Any line through a triangle that splits both the triangle's area and its perimeter in half goes through the triangle's incenter. There can be one, two, or three of these for any given triangle.
475
+
476
+ The formulas in this section are true for all Euclidean triangles.
477
+
478
+ The medians and the sides are related by[26]:p.70
479
+
480
+ and
481
+
482
+ and equivalently for mb and mc.
483
+
484
+ For angle A opposite side a, the length of the internal angle bisector is given by[27]
485
+
486
+ for semiperimeter s, where the bisector length is measured from the vertex to where it meets the opposite side.
487
+
488
+ The interior perpendicular bisectors are given by
489
+
490
+ where the sides are
491
+
492
+
493
+
494
+ a
495
+ ≥
496
+ b
497
+ ≥
498
+ c
499
+
500
+
501
+ {\displaystyle a\geq b\geq c}
502
+
503
+ and the area is
504
+
505
+
506
+
507
+ T
508
+ .
509
+
510
+
511
+ {\displaystyle T.}
512
+
513
+ [28]:Thm 2
514
+
515
+ The altitude from, for example, the side of length a is
516
+
517
+ The following formulas involve the circumradius R and the inradius r:
518
+
519
+ where ha etc. are the altitudes to the subscripted sides;[26]:p.79
520
+
521
+ and
522
+
523
+ The product of two sides of a triangle equals the altitude to the third side times the diameter D of the circumcircle:[26]:p.64
524
+
525
+ Suppose two adjacent but non-overlapping triangles share the same side of length f and share the same circumcircle, so that the side of length f is a chord of the circumcircle and the triangles have side lengths (a, b, f) and (c, d, f), with the two triangles together forming a cyclic quadrilateral with side lengths in sequence (a, b, c, d). Then[29]:84
526
+
527
+ Let G be the centroid of a triangle with vertices A, B, and C, and let P be any interior point. Then the distances between the points are related by[29]:174
528
+
529
+ The sum of the squares of the triangle's sides equals three times the sum of the squared distances of the centroid from the vertices:
530
+
531
+ Let qa, qb, and qc be the distances from the centroid to the sides of lengths a, b, and c. Then[29]:173
532
+
533
+ and
534
+
535
+ for area T.
536
+
537
+ Carnot's theorem states that the sum of the distances from the circumcenter to the three sides equals the sum of the circumradius and the inradius.[26]:p.83 Here a segment's length is considered to be negative if and only if the segment lies entirely outside the triangle. This method is especially useful for deducing the properties of more abstract forms of triangles, such as the ones induced by Lie algebras, that otherwise have the same properties as usual triangles.
538
+
539
+ Euler's theorem states that the distance d between the circumcenter and the incenter is given by[26]:p.85
540
+
541
+ or equivalently
542
+
543
+ where R is the circumradius and r is the inradius. Thus for all triangles R ≥ 2r, with equality holding for equilateral triangles.
544
+
545
+ If we denote that the orthocenter divides one altitude into segments of lengths u and v, another altitude into segment lengths w and x, and the third altitude into segment lengths y and z, then uv = wx = yz.[26]:p.94
546
+
547
+ The distance from a side to the circumcenter equals half the distance from the opposite vertex to the orthocenter.[26]:p.99
548
+
549
+ The sum of the squares of the distances from the vertices to the orthocenter H plus the sum of the squares of the sides equals twelve times the square of the circumradius:[26]:p.102
550
+
551
+ In addition to the law of sines, the law of cosines, the law of tangents, and the trigonometric existence conditions given earlier, for any triangle
552
+
553
+ Morley's trisector theorem states that in any triangle, the three points of intersection of the adjacent angle trisectors form an equilateral triangle, called the Morley triangle.
554
+
555
+ As discussed above, every triangle has a unique inscribed circle (incircle) that is interior to the triangle and tangent to all three sides.
556
+
557
+ Every triangle has a unique Steiner inellipse which is interior to the triangle and tangent at the midpoints of the sides. Marden's theorem shows how to find the foci of this ellipse.[31] This ellipse has the greatest area of any ellipse tangent to all three sides of the triangle.
558
+
559
+ The Mandart inellipse of a triangle is the ellipse inscribed within the triangle tangent to its sides at the contact points of its excircles.
560
+
561
+ For any ellipse inscribed in a triangle ABC, let the foci be P and Q. Then[32]
562
+
563
+ Every convex polygon with area T can be inscribed in a triangle of area at most equal to 2T. Equality holds (exclusively) for a parallelogram.[33]
564
+
565
+ The Lemoine hexagon is a cyclic hexagon with vertices given by the six intersections of the sides of a triangle with the three lines that are parallel to the sides and that pass through its symmedian point. In either its simple form or its self-intersecting form, the Lemoine hexagon is interior to the triangle with two vertices on each side of the triangle.
566
+
567
+ Every acute triangle has three inscribed squares (squares in its interior such that all four of a square's vertices lie on a side of the triangle, so two of them lie on the same side and hence one side of the square coincides with part of a side of the triangle). In a right triangle two of the squares coincide and have a vertex at the triangle's right angle, so a right triangle has only two distinct inscribed squares. An obtuse triangle has only one inscribed square, with a side coinciding with part of the triangle's longest side. Within a given triangle, a longer common side is associated with a smaller inscribed square. If an inscribed square has side of length qa and the triangle has a side of length a, part of which side coincides with a side of the square, then qa, a, the altitude ha from the side a, and the triangle's area T are related according to[34][35]
568
+
569
+ The largest possible ratio of the area of the inscribed square to the area of the triangle is 1/2, which occurs when a2 = 2T, q = a/2, and the altitude of the triangle from the base of length a is equal to a. The smallest possible ratio of the side of one inscribed square to the side of another in the same non-obtuse triangle is
570
+
571
+
572
+
573
+ 2
574
+
575
+
576
+ 2
577
+
578
+
579
+
580
+ /
581
+
582
+ 3
583
+ =
584
+ 0.94....
585
+
586
+
587
+ {\displaystyle 2{\sqrt {2}}/3=0.94....}
588
+
589
+ [35] Both of these extreme cases occur for the isosceles right triangle.
590
+
591
+ From an interior point in a reference triangle, the nearest points on the three sides serve as the vertices of the pedal triangle of that point. If the interior point is the circumcenter of the reference triangle, the vertices of the pedal triangle are the midpoints of the reference triangle's sides, and so the pedal triangle is called the midpoint triangle or medial triangle. The midpoint triangle subdivides the reference triangle into four congruent triangles which are similar to the reference triangle.
592
+
593
+ The Gergonne triangle or intouch triangle of a reference triangle has its vertices at the three points of tangency of the reference triangle's sides with its incircle. The extouch triangle of a reference triangle has its vertices at the points of tangency of the reference triangle's excircles with its sides (not extended).
594
+
595
+ The tangential triangle of a reference triangle (other than a right triangle) is the triangle whose sides are on the tangent lines to the reference triangle's circumcircle at its vertices.
596
+
597
+ As mentioned above, every triangle has a unique circumcircle, a circle passing through all three vertices, whose center is the intersection of the perpendicular bisectors of the triangle's sides.
598
+
599
+ Further, every triangle has a unique Steiner circumellipse, which passes through the triangle's vertices and has its center at the triangle's centroid. Of all ellipses going through the triangle's vertices, it has the smallest area.
600
+
601
+ The Kiepert hyperbola is the unique conic which passes through the triangle's three vertices, its centroid, and its circumcenter.
602
+
603
+ Of all triangles contained in a given convex polygon, there exists a triangle with maximal area whose vertices are all vertices of the given polygon.[36]
604
+
605
+ One way to identify locations of points in (or outside) a triangle is to place the triangle in an arbitrary location and orientation in the Cartesian plane, and to use Cartesian coordinates. While convenient for many purposes, this approach has the disadvantage of all points' coordinate values being dependent on the arbitrary placement in the plane.
606
+
607
+ Two systems avoid that feature, so that the coordinates of a point are not affected by moving the triangle, rotating it, or reflecting it as in a mirror, any of which give a congruent triangle, or even by rescaling it to give a similar triangle:
608
+
609
+ A non-planar triangle is a triangle which is not contained in a (flat) plane. Some examples of non-planar triangles in non-Euclidean geometries are spherical triangles in spherical geometry and hyperbolic triangles in hyperbolic geometry.
610
+
611
+ While the measures of the internal angles in planar triangles always sum to 180°, a hyperbolic triangle has measures of angles that sum to less than 180°, and a spherical triangle has measures of angles that sum to more than 180°. A hyperbolic triangle can be obtained by drawing on a negatively curved surface, such as a saddle surface, and a spherical triangle can be obtained by drawing on a positively curved surface such as a sphere. Thus, if one draws a giant triangle on the surface of the Earth, one will find that the sum of the measures of its angles is greater than 180°; in fact it will be between 180° and 540°.[37] In particular it is possible to draw a triangle on a sphere such that the measure of each of its internal angles is equal to 90°, adding up to a total of 270°.
612
+
613
+ Specifically, on a sphere the sum of the angles of a triangle is
614
+
615
+ where f is the fraction of the sphere's area which is enclosed by the triangle. For example, suppose that we draw a triangle on the Earth's surface with vertices at the North Pole, at a point on the equator at 0° longitude, and a point on the equator at 90° West longitude. The great circle line between the latter two points is the equator, and the great circle line between either of those points and the North Pole is a line of longitude; so there are right angles at the two points on the equator. Moreover, the angle at the North Pole is also 90° because the other two vertices differ by 90° of longitude. So the sum of the angles in this triangle is 90° + 90° + 90° = 270°. The triangle encloses 1/4 of the northern hemisphere (90°/360° as viewed from the North Pole) and therefore 1/8 of the Earth's surface, so in the formula f = 1/8; thus the formula correctly gives the sum of the triangle's angles as 270°.
616
+
617
+ From the above angle sum formula we can also see that the Earth's surface is locally flat: If we draw an arbitrarily small triangle in the neighborhood of one point on the Earth's surface, the fraction f of the Earth's surface which is enclosed by the triangle will be arbitrarily close to zero. In this case the angle sum formula simplifies to 180°, which we know is what Euclidean geometry tells us for triangles on a flat surface.
618
+
619
+ Rectangles have been the most popular and common geometric form for buildings since the shape is easy to stack and organize; as a standard, it is easy to design furniture and fixtures to fit inside rectangularly shaped buildings. But triangles, while more difficult to use conceptually, provide a great deal of strength. As computer technology helps architects design creative new buildings, triangular shapes are becoming increasingly prevalent as parts of buildings and as the primary shape for some types of skyscrapers as well as building materials. In Tokyo in 1989, architects had wondered whether it was possible to build a 500-story tower to provide affordable office space for this densely packed city, but with the danger to buildings from earthquakes, architects considered that a triangular shape would be necessary if such a building were to be built.[38]
620
+
621
+ In New York City, as Broadway crisscrosses major avenues, the resulting blocks are cut like triangles, and buildings have been built on these shapes; one such building is the triangularly shaped Flatiron Building which real estate people admit has a "warren of awkward spaces that do not easily accommodate modern office furniture" but that has not prevented the structure from becoming a landmark icon.[39] Designers have made houses in Norway using triangular themes.[40] Triangle shapes have appeared in churches[41] as well as public buildings including colleges[42] as well as supports for innovative home designs.[43]
622
+
623
+ Triangles are sturdy; while a rectangle can collapse into a parallelogram from pressure to one of its points, triangles have a natural strength which supports structures against lateral pressures. A triangle will not change shape unless its sides are bent or extended or broken or if its joints break; in essence, each of the three sides supports the other two. A rectangle, in contrast, is more dependent on the strength of its joints in a structural sense. Some innovative designers have proposed making bricks not out of rectangles, but with triangular shapes which can be combined in three dimensions.[44] It is likely that triangles will be used increasingly in new ways as architecture increases in complexity. It is important to remember that triangles are strong in terms of rigidity, but while packed in a tessellating arrangement triangles are not as strong as hexagons under compression (hence the prevalence of hexagonal forms in nature). Tessellated triangles still maintain superior strength for cantilevering however, and this is the basis for one of the strongest man made structures, the tetrahedral truss.
624
+
625
+
626
+
627
+
628
+
en/5791.html.txt ADDED
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1
+
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+
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+ Knitting is a method by which yarn is manipulated to create a textile or fabric; it is used in many types of garments. Knitting may be done by hand or by machine.
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+
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+ Knitting creates stitches: loops of yarn in a row, either flat or in the round (tubular). There are usually many active stitches on the knitting needle at one time. Knitted fabric consists of a number of consecutive rows of connected loops that intermesh with the next and previous rows. As each row is formed, each newly created loop is pulled through one or more loops from the prior row and placed on the gaining needle so that the loops from the prior row can be pulled off the other needle without unraveling.
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+
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+ Differences in yarn (varying in fibre type, weight, uniformity and twist), needle size, and stitch type allow for a variety of knitted fabrics with different properties, including color, texture, thickness, heat retention, water resistance, and integrity. A small sample of knitwork is known as a swatch.
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+
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+ Like weaving, knitting is a technique for producing a two-dimensional fabric made from a one-dimensional yarn or thread. In weaving, threads are always straight, running parallel either lengthwise (warp threads) or crosswise (weft threads). By contrast, the yarn in knitted fabrics follows a meandering path (a course), forming symmetric loops (also called bights) symmetrically above and below the mean path of the yarn. These meandering loops can be easily stretched in different directions giving knit fabrics much more elasticity than woven fabrics. Depending on the yarn and knitting pattern, knitted garments can stretch as much as 500%. For this reason, knitting was initially developed for garments that must be elastic or stretch in response to the wearer's motions, such as socks and hosiery. For comparison, woven garments stretch mainly along one or other of a related pair of directions that lie roughly diagonally between the warp and the weft, while contracting in the other direction of the pair (stretching and contracting with the bias), and are not very elastic, unless they are woven from stretchable material such as spandex. Knitted garments are often more form-fitting than woven garments, since their elasticity allows them to contour to the body's outline more closely; by contrast, curvature is introduced into most woven garments only with sewn darts, flares, gussets and gores, the seams of which lower the elasticity of the woven fabric still further. Extra curvature can be introduced into knitted garments without seams, as in the heel of a sock; the effect of darts, flares, etc. can be obtained with short rows or by increasing or decreasing the number of stitches. Thread used in weaving is usually much finer than the yarn used in knitting, which can give the knitted fabric more bulk and less drape than a woven fabric.
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+ If they are not secured, the loops of a knitted course will come undone when their yarn is pulled; this is known as ripping out, unravelling knitting, or humorously, frogging (because you 'rip it', this sounds like a frog croaking: 'rib-bit').[1] To secure a stitch, at least one new loop is passed through it. Although the new stitch is itself unsecured ("active" or "live"), it secures the stitch(es) suspended from it. A sequence of stitches in which each stitch is suspended from the next is called a wale.[2] To secure the initial stitches of a knitted fabric, a method for casting on is used; to secure the final stitches in a wale, one uses a method of binding/casting off. During knitting, the active stitches are secured mechanically, either from individual hooks (in knitting machines) or from a knitting needle or frame in hand-knitting.
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+
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+ There are two major varieties of knitting: weft knitting and warp knitting.[3] In the more common weft knitting, the wales are perpendicular to the course of the yarn. In warp knitting, the wales and courses run roughly parallel. In weft knitting, the entire fabric may be produced from a single yarn, by adding stitches to each wale in turn, moving across the fabric as in a raster scan. By contrast, in warp knitting, one yarn is required for every wale. Since a typical piece of knitted fabric may have hundreds of wales, warp knitting is typically done by machine, whereas weft knitting is done by both hand and machine.[4] Warp-knitted fabrics such as tricot and milanese are resistant to runs, and are commonly used in lingerie.
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+ Weft-knit fabrics may also be knit with multiple yarns, usually to produce interesting color patterns. The two most common approaches are intarsia and stranded colorwork. In intarsia, the yarns are used in well-segregated regions, e.g., a red apple on a field of green; in that case, the yarns are kept on separate spools and only one is knitted at any time. In the more complex stranded approach, two or more yarns alternate repeatedly within one row and all the yarns must be carried along the row, as seen in Fair Isle sweaters. Double knitting can produce two separate knitted fabrics simultaneously (e.g., two socks). However, the two fabrics are usually integrated into one, giving it great warmth and excellent drape.
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+ In securing the previous stitch in a wale, the next stitch can pass through the previous loop from either below or above. If the former, the stitch is denoted as a 'knit stitch' or a 'plain stitch;' if the latter, as a 'purl stitch'. The two stitches are related in that a knit stitch seen from one side of the fabric appears as a purl stitch on the other side.
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+ The two types of stitches have a different visual effect; the knit stitches look like 'V's stacked vertically, whereas the purl stitches look like a wavy horizontal line across the fabric. Patterns and pictures can be created in knitted fabrics by using knit and purl stitches as "pixels"; however, such pixels are usually rectangular, rather than square, depending on the gauge/tension of the knitting. Individual stitches, or rows of stitches, may be made taller by drawing more yarn into the new loop (an elongated stitch), which is the basis for uneven knitting: a row of tall stitches may alternate with one or more rows of short stitches for an interesting visual effect. Short and tall stitches may also alternate within a row, forming a fish-like oval pattern.
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+
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+ In the simplest of hand-knitted fabrics, every row of stitches are all knit (or all purl); this creates a garter stitch fabric. Alternating rows of all knit stitches and all purl stitches creates a stockinette pattern/stocking stitch. Vertical stripes (ribbing) are possible by having alternating wales of knit and purl stitches. For example, a common choice is 2x2 ribbing, in which two wales of knit stitches are followed by two wales of purl stitches, etc. Horizontal striping (welting) is also possible, by alternating rows of knit and purl stitches. Checkerboard patterns (basketweave) are also possible, the smallest of which is known as seed/moss stitch: the stitches alternate between knit and purl in every wale and along every row.
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+ Fabrics in which each knitted row is followed by a purled row, such as in stockinette/stocking stitch, have a tendency to curl—top and bottom curl toward the front (or knitted side) while the sides curl toward the back (or purled side); by contrast, those in which knit and purl stitches are arranged symmetrically (such as ribbing, garter stitch or seed/moss stitch) have more texture and tend to lie flat. Wales of purl stitches have a tendency to recede, whereas those of knit stitches tend to come forward, giving the fabric more stretchability. Thus, the purl wales in ribbing tend to be invisible, since the neighboring knit wales come forward. Conversely, rows of purl stitches tend to form an embossed ridge relative to a row of knit stitches. This is the basis of shadow knitting, in which the appearance of a knitted fabric changes when viewed from different directions.[5]
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+ Typically, a new stitch is passed through a single unsecured ('active') loop, thus lengthening that wale by one stitch. However, this need not be so; the new loop may be passed through an already secured stitch lower down on the fabric, or even between secured stitches (a dip stitch). Depending on the distance between where the loop is drawn through the fabric and where it is knitted, dip stitches can produce a subtle stippling or long lines across the surface of the fabric, e.g., the lower leaves of a flower. The new loop may also be passed between two stitches in the 'present' row, thus clustering the intervening stitches; this approach is often used to produce a smocking effect in the fabric. The new loop may also be passed through 'two or more' previous stitches, producing a decrease and merging wales together. The merged stitches need not be from the same row; for example, a tuck can be formed by knitting stitches together from two different rows, producing a raised horizontal welt on the fabric.
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+ Not every stitch in a row need be knitted; some may be 'missed' (unknitted and passed to the active needle) and knitted on a subsequent row. This is known as slip-stitch knitting.[6] The slipped stitches are naturally longer than the knitted ones. For example, a stitch slipped for one row before knitting would be roughly twice as tall as its knitted counterparts. This can produce interesting visual effects, although the resulting fabric is more rigid because the slipped stitch 'pulls' on its neighbours and is less deformable. Mosaic knitting is a form of slip-stitch knitting that knits alternate colored rows and uses slip stitches to form patterns; mosaic-knit fabrics tend to be stiffer than patterned fabrics produced by other methods such as Fair-Isle knitting.[7]
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+ In some cases, a stitch may be deliberately left unsecured by a new stitch and its wale allowed to disassemble. This is known as drop-stitch knitting, and produces a vertical ladder of see-through holes in the fabric, corresponding to where the wale had been.
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+ Both knit and purl stitches may be twisted: usually once if at all, but sometimes twice and (very rarely) thrice. When seen from above, the twist can be clockwise (right yarn over left) or counterclockwise (left yarn over right); these are denoted as right- and left-plaited stitches, respectively. Hand-knitters generally produce right-plaited stitches by knitting or purling through the back loops, i.e., passing the needle through the initial stitch in an unusual way, but wrapping the yarn as usual. By contrast, the left-plaited stitch is generally formed by hand-knitters by wrapping the yarn in the opposite way, rather than by any change in the needle. Although they are mirror images in form, right- and left-plaited stitches are functionally equivalent. Both types of plaited stitches give a subtle but interesting visual texture, and tend to draw the fabric inwards, making it stiffer. Plaited stitches are a common method for knitting jewelry from fine metal wire.
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+ The initial and final edges of a knitted fabric are known as the cast-on and bound/cast-off edges. The side edges are known as the selvages; the word derives from "self-edges", meaning that the stitches do not need to be secured by anything else. Many types of selvages have been developed, with different elastic and ornamental properties.
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+ Vertical and horizontal edges can be introduced within a knitted fabric, e.g., for button holes, by binding/casting off and re-casting on again (horizontal) or by knitting the fabrics on either side of a vertical edge separately.
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+ Two knitted fabrics can be joined by embroidery-based grafting methods, most commonly the Kitchener stitch. New wales can be begun from any of the edges of a knitted fabric; this is known as picking up stitches and is the basis for entrelac, in which the wales run perpendicular to one another in a checkerboard pattern.
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+ Ordinarily, stitches are knitted in the same order in every row, and the wales of the fabric run parallel and vertically along the fabric. However, this need not be so, since the order in which stitches are knitted may be permuted so that wales cross over one another, forming a cable pattern. Cables patterns tend to draw the fabric together, making it denser and less elastic;[8] Aran sweaters are a common form of knitted cabling.[9] Arbitrarily complex braid patterns can be done in cable knitting, with the proviso that the wales must move ever upwards; it is generally impossible for a wale to move up and then down the fabric. Knitters have developed methods for giving the illusion of a circular wale, such as appear in Celtic knots, but these are inexact approximations. However, such circular wales are possible using Swiss darning, a form of embroidery, or by knitting a tube separately and attaching it to the knitted fabric.
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+ A wale can split into two or more wales using increases, most commonly involving a yarn over. Depending on how the increase is done, there is often a hole in the fabric at the point of the increase. This is used to great effect in lace knitting, which consists of making patterns and pictures using such holes, rather than with the stitches themselves.[10] The large and many holes in lacy knitting makes it extremely elastic; for example, some Shetland "wedding-ring" shawls are so fine that they may be drawn through a wedding ring.
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+ By combining increases and decreases, it is possible to make the direction of a wale slant away from vertical, even in weft knitting. This is the basis for bias knitting, and can be used for visual effect, similar to the direction of a brush-stroke in oil painting.
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+
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+ Various point-like ornaments may be added to knitting for their look or to improve the wear of the fabric. Examples include various types of bobbles, sequins and beads. Long loops can also be drawn out and secured, forming a "shaggy" texture to the fabric; this is known as loop knitting. Additional patterns can be made on the surface of the knitted fabric using embroidery; if the embroidery resembles knitting, it is often called Swiss darning. Various closures for the garments, such as frogs and buttons can be added; usually buttonholes are knitted into the garment, rather than cut.
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+ Ornamental pieces may also be knitted separately and then attached using applique. For example, differently colored leaves and petals of a flower could be knit separately and attached to form the final picture. Separately knitted tubes can be applied to a knitted fabric to form complex Celtic knots and other patterns that would be difficult to knit.
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+ Unknitted yarns may be worked into knitted fabrics for warmth, as is done in tufting and "weaving" (also known as "couching").
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+ The word is derived from knot and ultimately from the Old English cnyttan, to knot.[11]
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+ The exact origins of knitting are unknown, the earliest known examples being cotton socks found in Egyptian pyramids.[12]
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+ Nålebinding (Danish: literally "binding with a needle" or "needle-binding") is a fabric creation technique predating both knitting and crochet.
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+ The first commercial knitting guilds appear in Western Europe in the early fifteenth century (Tournai in 1429, Barcelona in 1496). The Guild of Saint Fiacre was founded in Paris in 1527 but the archives mention an organization (not necessarily a guild) of knitters from 1268.[13]
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+ With the invention of the stocking frame, an early form of knitting machine, knitting "by hand" became a craft used by country people with easy access to fiber. Similar to quilting, spinning, and needlepoint, hand knitting became a leisure activity for the wealthy. English Roman Catholic priest and a former Anglican bishop, Richard Rutt, authored a history of the craft in A History of Hand Knitting (Batsford, 1987). His collection of books about knitting is now housed at the Winchester School of Art (University of Southampton).
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+ The topology of a knitted fabric is relatively complex. Unlike woven fabrics, where strands usually run straight horizontally and vertically, yarn that has been knitted follows a looped path along its row, as with the red strand in the diagram at left, in which the loops of one row have all been pulled through the loops of the row below it.
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+ Because there is no single straight line of yarn anywhere in the pattern, a knitted piece of fabric can stretch in all directions. This elasticity is all but unavailable in woven fabrics which only stretch along the bias. Many modern stretchy garments, even as they rely on elastic synthetic materials for some stretch, also achieve at least some of their stretch through knitted patterns.
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+ The basic knitted fabric (as in the diagram, and usually called a stocking or stockinette pattern) has a definite "right side" and "wrong side". On the right side, the visible portions of the loops are the verticals connecting two rows which are arranged in a grid of V shapes. On the wrong side, the ends of the loops are visible, both the tops and bottoms, creating a much more bumpy texture sometimes called reverse stockinette. (Despite being the "wrong side," reverse stockinette is frequently used as a pattern in its own right.) Because the yarn holding rows together is all on the front, and the yarn holding side-by-side stitches together is all on the back, stockinette fabric has a strong tendency to curl toward the front on the top and bottom, and toward the back on the left and right side.
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+ Stitches can be worked from either side, and various patterns are created by mixing regular knit stitches with the "wrong side" stitches, known as purl stitches, either in columns (ribbing), rows (garter, welting), or more complex patterns. Each fabric has different properties: a garter stitch has much more vertical stretch, while ribbing stretches much more horizontally. Because of their front-back symmetry, these two fabrics have little curl, making them popular as edging, even when their stretch properties are not desired.
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+ Different combinations of knit and purl stitches, along with more advanced techniques, generate fabrics of considerably variable consistency, from gauzy to very dense, from highly stretchy to relatively stiff, from flat to tightly curled, and so on.
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+ The most common texture for a knitted garment is that generated by the flat stockinette stitch—as seen, though very small, in machine-made stockings and T-shirts—which is worked in the round as nothing but knit stitches, and worked flat as alternating rows of knit and purl. Other simple textures can be made with nothing but knit and purl stitches, including garter stitch, ribbing, and moss and seed stitches. Adding a "slip stitch" (where a loop is passed from one needle to the other) allows for a wide range of textures, including heel and linen stitches as well as a number of more complicated patterns.
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+ Some more advanced knitting techniques create a surprising variety of complex textures. Combining certain increases, which can create small eyelet holes in the resulting fabric, with assorted decreases is key to creating knitted lace, a very open fabric resembling needle or bobbin lace. Open vertical stripes can be created using the drop-stitch knitting technique. Changing the order of stitches from one row to the next, usually with the help of a cable needle or stitch holder, is key to cable knitting, producing an endless variety of cables, honeycombs, ropes, and Aran sweater patterning. Entrelac forms a rich checkerboard texture by knitting small squares, picking up their side edges, and knitting more squares to continue the piece.
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+ Fair Isle knitting uses two or more colored yarns to create patterns and forms a thicker and less flexible fabric.
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+ The appearance of a garment is also affected by the weight of the yarn, which describes the thickness of the spun fibre. The thicker the yarn, the more visible and apparent stitches will be; the thinner the yarn, the finer the texture.
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+ Plenty of finished knitting projects never use more than a single color of yarn, but there are many ways to work in multiple colors. Some yarns are dyed to be either variegated (changing color every few stitches in a random fashion) or self-striping (changing every few rows). More complicated techniques permit large fields of color (intarsia, for example), busy small-scale patterns of color (such as Fair Isle), or both (double knitting and slip-stitch color, for example).
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+ Yarn with multiple shades of the same hue are called ombre, while a yarn with multiple hues may be known as a given colorway; a green, red and yellow yarn might be dubbed the "Parrot Colorway" by its manufacturer, for example. Heathered yarns contain small amounts of fibre of different colours, while tweed yarns may have greater amounts of different colored fibres.
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+ There are many hundreds of different knitting stitches used by hand knitters. A piece of hand knitting begins with the process of casting on, which involves the initial creation of the stitches on the needle. Different methods of casting on are used for different effects: one may be stretchy enough for lace, while another provides a decorative edging. Provisional cast-ons are used when the knitting will continue in both directions from the cast-on. There are various methods employed to cast on, such as the "thumb method" (also known as "slingshot" or "long-tail" cast-ons), where the stitches are created by a series of loops that will, when knitted, give a very loose edge ideal for "picking up stitches" and knitting a border; the "double needle method" (also known as "knit-on" or "cable cast-on"), whereby each loop placed on the needle is then "knitted on," which produces a firmer edge ideal on its own as a border; and many more. The number of active stitches remains the same as when cast on unless stitches are added (an increase) or removed (a decrease).
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+ Most Western-style hand knitters follow either the English style (in which the yarn is held in the right hand) or the Continental style (in which the yarn is held in the left hand).
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+ There are also different ways to insert the needle into the stitch. Knitting through the front of a stitch is called Western knitting. Going through the back of a stitch is called Eastern knitting. A third method, called combination knitting, goes through the front of a knit stitch and the back of a purl stitch.[14]
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+ Once the hand knitted piece is finished, the remaining live stitches are "cast off". Casting (or "binding") off loops the stitches across each other so they can be removed from the needle without unravelling the item. Although the mechanics are different from casting on, there is a similar variety of methods.
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+ In hand knitting certain articles of clothing, especially larger ones like sweaters, the final knitted garment will be made of several knitted pieces, with individual sections of the garment hand knitted separately and then sewn together. Seamless knitting, where a whole garment is hand knit as a single piece, is also possible. Elizabeth Zimmermann is probably the best-known proponent of seamless or circular hand knitting techniques. Smaller items, such as socks and hats, are usually knit in one piece on double-pointed needles or circular needles. Hats in particular can be started "top down" on double pointed needles with the increases added until the preferred size is achieved, switching to an appropriate circular needle when enough stitches have been added. Care must be taken to bind off at a tension that will allow the "give" needed to comfortably fit on the head. (See Circular knitting.)
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+ Yarn for hand-knitting is usually sold as balls or skeins (hanks), and it may also be wound on spools or cones. Skeins and balls are generally sold with a yarn-band, a label that describes the yarn's weight, length, dye lot, fiber content, washing instructions, suggested needle size, likely gauge/tension, etc. It is common practice to save the yarn band for future reference, especially if additional skeins must be purchased. Knitters generally ensure that the yarn for a project comes from a single dye lot. The dye lot specifies a group of skeins that were dyed together and thus have precisely the same color; skeins from different dye-lots, even if very similar in color, are usually slightly different and may produce a visible horizontal stripe when knitted together. If a knitter buys insufficient yarn of a single dye lot to complete a project, additional skeins of the same dye lot can sometimes be obtained from other yarn stores or online. Otherwise, knitters can alternate skeins every few rows to help the dye lots blend together easier.
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+ The thickness or weight of the yarn is a significant factor in determining the gauge/tension, i.e., how many stitches and rows are required to cover a given area for a given stitch pattern. Thicker yarns generally require thicker knitting needles, whereas thinner yarns may be knit with thick or thin needles. Hence, thicker yarns generally require fewer stitches, and therefore less time, to knit up a given garment. Patterns and motifs are coarser with thicker yarns; thicker yarns produce bold visual effects, whereas thinner yarns are best for refined patterns. Yarns are grouped by thickness into six categories: superfine, fine, light, medium, bulky and superbulky; quantitatively, thickness is measured by the number of wraps per inch (WPI). In the British Commonwealth (outside North America) yarns are measured as 1ply, 2ply, 3ply, 4ply, 5ply, 8ply (or double knit),10ply and 12ply (triple knit). The related weight per unit length is usually measured in tex or denier.
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+ Before knitting, the knitter will typically transform a hank/skein into a ball where the yarn emerges from the center of the ball; this making the knitting easier by preventing the yarn from becoming easily tangled. This transformation may be done by hand, or with a device known as a ballwinder. When knitting, some knitters enclose their balls in jars to keep them clean and untangled with other yarns; the free yarn passes through a small hole in the jar-lid.
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+ A yarn's usefulness for a knitting project is judged by several factors, such as its loft (its ability to trap air), its resilience (elasticity under tension), its washability and colorfastness, its hand (its feel, particularly softness vs. scratchiness), its durability against abrasion, its resistance to pilling, its hairiness (fuzziness), its tendency to twist or untwist, its overall weight and drape, its blocking and felting qualities, its comfort (breathability, moisture absorption, wicking properties) and of course its look, which includes its color, sheen, smoothness and ornamental features. Other factors include allergenicity; speed of drying; resistance to chemicals, moths, and mildew; melting point and flammability; retention of static electricity; and the propensity to become stained and to accept dyes. Different factors may be more significant than others for different knitting projects, so there is no one "best" yarn. The resilience and propensity to (un)twist are general properties that affect the ease of hand-knitting. More resilient yarns are more forgiving of irregularities in tension; highly twisted yarns are sometimes difficult to knit, whereas untwisting yarns can lead to split stitches, in which not all the yarn is knitted into a stitch. A key factor in knitting is stitch definition, corresponding to how well complicated stitch patterns can be seen when made from a given yarn. Smooth, highly spun yarns are best for showing off stitch patterns; at the other extreme, very fuzzy yarns or eyelash yarns have poor stitch definition, and any complicated stitch pattern would be invisible.
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+
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+ Although knitting may be done with ribbons, metal wire or more exotic filaments, most yarns are made by spinning fibers. In spinning, the fibers are twisted so that the yarn resists breaking under tension; the twisting may be done in either direction, resulting in a Z-twist or S-twist yarn. If the fibers are first aligned by combing them, the yarn is smoother and called a worsted; by contrast, if the fibers are carded but not combed, the yarn is fuzzier and called woolen-spun. The fibers making up a yarn may be continuous filament fibers such as silk and many synthetics, or they may be staples (fibers of an average length, typically a few inches); naturally filament fibers are sometimes cut up into staples before spinning. The strength of the spun yarn against breaking is determined by the amount of twist, the length of the fibers and the thickness of the yarn. In general, yarns become stronger with more twist (also called worst), longer fibers and thicker yarns (more fibers); for example, thinner yarns require more twist than do thicker yarns to resist breaking under tension. The thickness of the yarn may vary along its length; a slub is a much thicker section in which a mass of fibers is incorporated into the yarn.
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+
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+ The spun fibers are generally divided into animal fibers, plant and synthetic fibers. These fiber types are chemically different, corresponding to proteins, carbohydrates and synthetic polymers, respectively. Animal fibers include silk, but generally are long hairs of animals such as sheep (wool), goat (angora, or cashmere goat), rabbit (angora), llama, alpaca, dog, cat, camel, yak, and muskox (qiviut). Plants used for fibers include cotton, flax (for linen), bamboo, ramie, hemp, jute, nettle, raffia, yucca, coconut husk, banana fiber, soy and corn. Rayon and acetate fibers are also produced from cellulose mainly derived from trees. Common synthetic fibers include acrylics,[15] polyesters such as dacron and ingeo, nylon and other polyamides, and olefins such as polypropylene. Of these types, wool is generally favored for knitting, chiefly owing to its superior elasticity, warmth and (sometimes) felting. It is also common to blend different fibers in the yarn, e.g., 85% alpaca and 15% silk. Even within a type of fiber, there can be great variety in the length and thickness of the fibers; for example, Merino wool and Egyptian cotton are favored because they produce exceptionally long, thin (fine) fibers for their type.
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+
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+ A single spun yarn may be knitted as is, or braided or plied with another. In plying, two or more yarns are spun together, almost always in the opposite sense from which they were spun individually; for example, two Z-twist yarns are usually plied with an S-twist. The opposing twist relieves some of the yarns' tendency to curl up and produces a thicker, balanced yarn. Plied yarns may themselves be plied together, producing cabled yarns or multi-stranded yarns. Sometimes, the yarns being plied are fed at different rates, so that one yarn loops around the other, as in bouclé. The single yarns may be dyed separately before plying, or afterwards to give the yarn a uniform look.
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+
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+ The dyeing of yarns is a complex art that has a long history. However, yarns need not be dyed. They may be dyed just one color, or a great variety of colors. Dyeing may be done industrially, by hand or even hand-painted onto the yarn. A great variety of synthetic dyes have been developed since the synthesis of indigo dye in the mid-19th century; however, natural dyes are also possible, although they are generally less brilliant. The color-scheme of a yarn is sometimes called its colorway. Variegated yarns can produce interesting visual effects, such as diagonal stripes; conversely, a variegated yarn may obscure a detailed knitting design, such as a cable or lace pattern.
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+ There are multiple commercial applications for knit fabric made of metal wire by knitting machines. Steel wire of various sizes may be used for electric and magnetic shielding due to its conductivity. Stainless steel may be used in a coffee press for its rust resistance.
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+ Metal wire can also be used as jewelry.
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+ Knitted glass combines knitting, lost-wax casting, mold-making, and kiln-casting.
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+ The process involves
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+ The process of knitting has three basic tasks:
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+ In very simple cases, knitting can be done without tools, using only the fingers to do these tasks; however, knitting is usually carried out using tools such as knitting needles, knitting machines or rigid frames. Depending on their size and shape, the rigid frames are called stocking frames, knitting boards, knitting rings (also called knitting looms) or knitting spools (also known as knitting knobbies, knitting nancies, or corkers). There is also a technique called knooking[18] of knitting with a crochet hook that has a cord attached to the end, to hold the stitches while they're being worked. Other tools are used to prepare yarn for knitting, to measure and design knitted garments, or to make knitting easier or more comfortable.
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+
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+ There are three basic types of knitting needles (also called "knitting pins"). The first and most common type consists of two slender, straight sticks tapered to a point at one end, and with a knob at the other end to prevent stitches from slipping off. Such needles are usually 10–16 inches (250–410 mm) long but, due to the compressibility of knitted fabrics, may be used to knit pieces significantly wider. The most important property of needles is their diameter, which ranges from below 2 to 25 mm (roughly 1 inch). The diameter affects the size of stitches, which affects the gauge/tension of the knitting and the elasticity of the fabric. Thus, a simple way to change gauge/tension is to use different needles, which is the basis of uneven knitting. Although the diameter of the knitting needle is often measured in millimeters, there are several measurement systems, particularly those specific to the United States, the United Kingdom and Japan; a conversion table is given at knitting needle. Such knitting needles may be made out of any materials, but the most common materials are metals, wood, bamboo, and plastic. Different materials have different frictions and grip the yarn differently; slick needles such as metallic needles are useful for swift knitting, whereas rougher needles such as bamboo offer more friction and are therefore less prone to dropping stitches. The knitting of new stitches occurs only at the tapered ends. Needles with lighted tips have been sold to allow knitters to knit in the dark.
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+
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+ The second type of knitting needles are straight, double-pointed knitting needles (also called "DPNs"). Double-pointed needles are tapered at both ends, which allows them to be knit from either end. DPNs are typically used for circular knitting, especially smaller tube-shaped pieces such as sleeves, collars, and socks; usually one needle is active while the others hold the remaining active stitches. DPNs are somewhat shorter (typically 7 inches) and are usually sold in sets of four or five.
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+
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+ The third needle type consists of circular needles, which are long, flexible double-pointed needles. The two tapered ends (typically 5 inches (130 mm) long) are rigid and straight, allowing for easy knitting; however, the two ends are connected by a flexible strand (usually nylon) that allows the two ends to be brought together. Circular needles are typically 24-60 inches long, and are usually used singly or in pairs; again, the width of the knitted piece may be significantly longer than the length of the circular needle. Interchangeable needles are a subset of circular needles. They are kits consist of pairs of needles with usually nylon cables or cords. The cables/cords are screwed into the needles, allowing the knitter to have both flexible straight needles or circular needles. This also allows the knitter to change the diameter and length of the needles as needed. The needles must be screwed on tightly, otherwise yarn can snag and become damaged.
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+
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+ The ability to work from either end of one needle is convenient in several types of knitting, such as slip-stitch versions of double knitting. Circular needles may be used for flat or circular knitting.
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+ Cable needles are a special case of DPNs, although they are usually not straight, but dimpled in the middle. Often, they have the form of a hook. When cabling a knitted piece, a hook is easier to grab and hold the yarn. Cable needles are typically very short (a few inches), and are used to hold stitches temporarily while others are being knitted. When in use, the cable needle is used at the same time as two regular needles. At specific points indicated by the knitting pattern, the cable needle is moved, the stitches on it are worked by the other needles, then the cable needle is turned around to a different position to create the cable twist.
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+
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+ Cable needles are a specific design, and are used to create the twisting motif of a knitted cable. They are made in different sizes, which produces cables of different widths.
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+ The largest aluminum circular knitting needles on record are size US 150 and are nearly 7 feet tall. They are owned by Paradise Fibers and are currently on display in the Paradise Fibers retail showroom.
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+ The current holder of the Guinness World Record for Knitting with the Largest Knitting Needles is Julia Hopson[19] of Penzance in Cornwall.
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+ Julia knitted a square of ten stitches and ten rows in stockinette stitch using knitting needles that were 6.5 centimeters in diameter and 3.5 meters long.
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+
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+ Various tools have been developed to make hand-knitting easier. Tools for measuring needle diameter and yarn properties have been discussed above, as well as the yarn swift, ballwinder and "yarntainers". Crochet hooks and a darning needle are often useful in binding/casting off or in joining two knitted pieces edge-to-edge. The darning needle is used in duplicate stitch (also known as Swiss darning). The crochet hook is also essential for repairing dropped stitches and some specialty stitches such as tufting. Other tools such as knitting spools or pom-pom makers are used to prepare specific ornaments. For large or complex knitting patterns, it is sometimes difficult to keep track of which stitch should be knit in a particular way; therefore, several tools have been developed to identify the number of a particular row or stitch, including circular stitch markers, hanging markers, extra yarn and row counters. A second potential difficulty is that the knitted piece will slide off the tapered end of the needles when unattended; this is prevented by "point protectors" that cap the tapered ends. Another problem is that too much knitting may lead to hand and wrist troubles; for this, special stress-relieving gloves are available. In traditional Shetland knitting a special belt is often used to support the end of one needle allowing the knitting greater speed. Finally, there are sundry bags and containers for holding knitting, yarns and needles.
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+
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+ Continental knitting is achieved by holding the yarn in your left hand for both knitting and purling. Patterns are created on the outside (public-facing) side of the piece.
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+ English-style knitting is achieved by holding the yarn in your right hand. Patterns are created on the outside (public-facing) side of the piece.
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+ This style is achieved by carrying the yarn around the neck or from a necklace-style hook, allowing the knitter to knit on the reverse (purl) side, e.g. "inside out" compared to Western knitting techniques. Patterns are typically created by stranding the yarn on the outside of the piece.
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+ Mega knitting is a term recently coined and relates to the use of knitting needles greater than or equal to half an inch in diameter.
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+ Mega knitting uses the same stitches and techniques as conventional knitting, except that hooks are carved into the ends of the needles. The hooked needles greatly enhance control of the work, catching the stitches and preventing them from slipping off.
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+ It was the development of the knitting machine that introduced hooked needles and enabled faultless, automated knitting. The hook catches the loop of yarn as each stitch is knitted, meaning that wrists and fingers do not have to work so hard and there is less chance of stitches slipping off the needle. The position of the hook is most important. Turn the left (non-working) hook to face away at all times; turn the right (working) hook toward you up whilst knitting (plain stitch) and away whilst purling.
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+ Mega knitting produces a chunky, bulky fabric or an open lacy weave, depending on the weight and type of yarn used.[20]
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+ Micro knitting or miniature knitting uses extremely fine threads and needles. Anthea Crome created 14 tiny sweaters used in the stop motion animated film Coraline and has made objects at 60 or 80 stitches per inch, making her own needles from fine surgical steel wire.[21][22][23] She has published Bugknits: Extreme knitting for hobbyists, artists and knitters (2009, Blurb: ISBN 978-1320025546). Annelies de Kort has knitted on an even smaller scale and has used needles of 0.4mm.[24][25]
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+ Industrially, metal wire is also knitted into a metal fabric for a wide range of uses including the filter material in cafetieres, catalytic converters for cars and many other uses. These fabrics are usually manufactured on circular knitting machines that would be recognized by conventional knitters as sock machines.
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+ Many fashion designers make heavy use of knitted fabric in their fashion collections. Gordana Gelhausen, who appeared in season six of the television show Project Runway, is primarily a knit designer. Other designers and labels that make heavy use of knitting include Michael Kors, Fendi, and Marc Jacobs.
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+ For individual hobbyists, websites such as Etsy, Big Cartel and Ravelry have made it easy to sell knitting patterns on a small scale, in a way similar to eBay.
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+ In the last decade, a practice called knitting graffiti, guerilla knitting, or yarn bombing—the use of knitted or crocheted cloth to modify and beautify one's (usually outdoor) surroundings—emerged in the U.S. and spread worldwide.[26] Magda Sayeg is credited with starting the movement in the US and Knit the City are a prominent group of graffiti knitters in the United Kingdom.[27] Yarn bombers sometimes target existing pieces of graffiti for beautification. For instance, Dave Cole is a contemporary sculpture artist who practiced knitting as graffiti for a large-scale public art installation in Melbourne, Australia for the Big West Arts Festival in 2009. The work was vandalized the night of its completion.[28] A new movie, shot by a Tasmanian filmmaker on a set made almost entirely out of yarn, was partially inspired by "knitted graffiti".[29]
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+ Many major metropolitan cities across the US and Europe host annual Yarn Crawls. The event is typically a multi-day event that caters to all knitters, crochet and yarn enthusiasts that supports the local crafting community. Over the multi-day period, multiple local yarn and knit shops participate in the yarn crawl and offer up store discounts, give away free exclusive patterns, provide classes, trunk shows and conduct raffles for prizes. Participants of the crawl receive a passport and get their passport stamped at each store they visit along the crawl. Traditionally those that get their passports fully stamped are eligible to win a larger gift basket filled with yarn, knitting and crochet goodies. Some local crawls also provide a Knit-Along (KAL) or Crochet-Along (CAL) where attendees follow a specific pattern prior to the crawl and then proudly wear it during the crawl for others to see.[30][31][32][33]
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+ Hand knitting garments for free distribution to others has become common practice among hand knitting groups. Girls and women hand knitted socks, sweaters, scarves, mittens, gloves, and hats for soldiers in Crimea, the American Civil War, and the Boer Wars; this practice continued in World War I, World War II and the Korean War, and continues for soldiers in Iraq and Afghanistan. The Australian charity Wrap with Love continues to provide blankets hand knitted by volunteers to people most in need around the world who have been affected by war.
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+ In the historical projects, yarn companies provided knitting patterns approved by the various branches of the armed services; often they were distributed by local chapters of the American Red Cross. Modern projects usually entail the hand knitting of hats or helmet liners; the liners provided for soldiers must be of 100% worsted weight wool and be crafted using specific colors.
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+ Clothing and afghans are frequently made for children, the elderly, and the economically disadvantaged in various countries. Pine Ridge Indian Reservation accepts donations for the Lakota people in the United States. Prayer shawls, or shawls in which the crafter meditates or says prayers of their faith while hand knitting with the intent on comforting the recipient, are donated to those experiencing loss or stress. Many knitters today hand knit and donate "chemo caps," soft caps for cancer patients who lose their hair during chemotherapy. Yarn companies offer free knitting patterns for these caps.
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+ Penguin sweaters were hand knitted by volunteers for the rehabilitation of penguins contaminated by exposure to oil slicks. The project is now complete.[34]
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+ Chicken sweaters were also hand knitted to aid battery hens that had lost their feathers. The organization is not currently accepting donations, but maintains a list of volunteers.[35]
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+ Originally started after the 2004 Indonesian tsunami, Knitters Without Borders[36] is a charity challenge issued by knitting personality Stephanie Pearl-McPhee that encourages hand knitters to donate to Médecins Sans Frontières (Doctors Without Borders). Instead of[hand knitting for charity, knitters are encouraged to donate a week's worth of disposable income, including money that otherwise might have been spent on yarn. Knitted items are occasional offered as prizes to donors. As of September 2011, Knitters Without Borders donors have contributed CAD$1,062,217.[37]
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+ Security blankets can also be made through the Project Linus organization which helps needy children.[38]
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+ There are organizations that help reach other countries in need such as afghans for Afghans. This outreach is described as, "afghans for Afghans is a humanitarian and educational people-to-people project that sends hand-knit and crocheted blankets and sweaters, vests, hats, mittens, and socks to the beleaguered people of Afghanistan."[39]
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+ The knitters of the Little Yellow Duck Project craft small yellow ducks which are left for others to find, as a random act of kindness and to raise awareness of blood donation and organ donation. The project was started in memory of a young woman who had collected plastic toy ducks and who died from cystic fibrosis while waiting for a lung transplant. Finders of the ducks are encouraged to log them on a website, which as of May 2020[update] shows that 12,265 ducks have been found in 106 countries.[40]
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+ Studies have shown that hand knitting, along with other forms of needlework, provide several significant health benefits. These studies have found the rhythmic and repetitive action of hand knitting can
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+ help prevent and manage stress, pain and depression, which in turn strengthens the body's immune system,[41] as well as create a relaxation response in the body which can decrease blood pressure, heart rate, help prevent illness, and have a calming effect. Pain specialists have also found that hand knitting changes brain chemistry, resulting in an increase in "feel good" hormones (i.e. serotonin and dopamine) and a decrease in stress hormones.[41]
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+ Hand knitting, along with other leisure activities, has been linked to reducing the risk of developing Alzheimer's disease and dementia. Much like physical activity strengthens the body, mental exercise makes the human brain more resilient.[42]
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+ A repository of research into the effect on health of hand knitting can be found at Stitch links,[43] an organization founded in Bath, England.
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+ Knitting also helps in the area of social interaction; knitting provides people with opportunities to socialize with others. Some ways to increase social interaction with knitting is inviting friends over to knit and chat with each other. Even if they've never knitted before this can be a fun way to interact with friends.[2] Many public libraries and yarn stores host knitting groups where knitters can meet locally to engage with others interested in hand crafts.
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+ Another interesting way that knitting can positively impact one's life is improving the dexterity in your hands and fingers. This keeps the fingers limber and can be especially helpful for those with arthritis. Knitting can reduce the pain of arthritis if people make it a daily habit.[2]
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+ Trinidad and Tobago (/ˈtrɪnɪdæd ... təˈbeɪɡoʊ/ (listen), /- toʊ-/), officially the Republic of Trinidad and Tobago, is the southernmost island country in the Caribbean.[14][15] Consisting of the main islands Trinidad and Tobago, and numerous much smaller islands, it is situated 130 kilometres (81 miles) south of Grenada and 11 kilometres (6.8 miles) off the coast of northeastern Venezuela.[16] It shares maritime boundaries with Barbados to the northeast, Grenada to the northwest, Guyana to the southeast, and Venezuela to the south and west.[17][18]
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+ The island of Trinidad was inhabited for centuries by native Amerindian peoples before becoming a colony in the Spanish Empire, following the arrival of Christopher Columbus in 1498. Spanish governor Don José María Chacón surrendered the island to a British fleet under the command of Sir Ralph Abercromby in 1797.[19] During the same period, the island of Tobago changed hands among Spanish, British, French, Dutch and Courlander colonisers more times than any other island in the Caribbean.[citation needed] Trinidad and Tobago were ceded to Britain in 1802 under the Treaty of Amiens as separate states and unified in 1889.[20] Trinidad and Tobago obtained independence in 1962, becoming a republic in 1976.[15][16]
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+ Trinidad and Tobago has the third highest GDP per capita based on purchasing power parity (PPP) in the Americas after the United States and Canada.[21] It is recognised by the World Bank as a high-income economy.[22] Unlike most Caribbean nations and territories, which rely heavily on tourism, the Trinidadian economy is primarily industrial with an emphasis on petroleum and petrochemicals;[23] much of the nation's wealth is derived from its large reserves of oil and natural gas.[24][25]
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+ Trinidad and Tobago is well known for its African and Indian cultures, reflected in its large and famous Carnival, Diwali, and Hosay celebrations, as well being the birthplace of steelpan, the limbo, and music styles such as calypso, soca, rapso, parang, chutney, and chutney soca.[26][27][28][29][30][31][32][33]
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+ Historian E. L. Joseph claimed that Trinidad's Amerindian name was Cairi or "Land of the Humming Bird", derived from the Arawak name for hummingbird, ierèttê or yerettê.[citation needed] However, other authors dispute this etymology with some claiming that cairi does not mean hummingbird (tukusi or tucuchi being suggested as the correct word) and some claiming that kairi, or iere, simply means island.[clarification needed][34][citation needed] Christopher Columbus renamed it "La Isla de la Trinidad" ("The Island of the Trinity"), fulfilling a vow made before setting out on his third voyage of exploration.[35] Tobago's cigar-like shape, or the use of tobacco by the native people, may have given it its Spanish name (cabaco, tavaco, tobacco) and possibly some of its other Amerindian names, such as Aloubaéra (black conch) and Urupaina (big snail),[34] although the English pronunciation is /təˈbeɪɡoʊ/.[citation needed]
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+ Both Trinidad and Tobago were originally settled by Amerindians who came through South America.[16] Trinidad was first settled by pre-agricultural Archaic people at least 7,000 years ago, making it the earliest settled part of the Caribbean.[36] Banwari Trace in south-west Trinidad is the oldest attested archaeological site in the Caribbean, dating to about 5000 BC. Several waves of migration occurred over the following centuries, which can be identified by differences in their archaeological remains.[37] At the time of European contact, Trinidad was occupied by various Arawakan-speaking groups including the Nepoya and Suppoya, and Cariban-speaking groups such as the Yao, while Tobago was occupied by the Island Caribs and Galibi. Trinidad was known to the native peoples as 'Ieri' ('Land of the Humming Bird').[36]
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+ Christopher Columbus was the first European to see Trinidad, on his third voyage to the Americas in 1498.[36][38] He also reported seeing Tobago on the distant horizon, naming it Bellaforma, but did not land on the island.[16][39]
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+ In the 1530s Antonio de Sedeño, a Spanish soldier intent on conquering the island of Trinidad, landed on its southwest coast with a small army of men, intending to subdue the Amerindian peoples of the island. Sedeño and his men fought the native peoples on many occasions, and subsequently built a fort. The next few decades were generally spent in warfare with the native peoples, until in 1592, the 'Cacique' (native chief) Wannawanare (also known as Guanaguanare) granted the area around modern Saint Joseph to Domingo de Vera e Ibargüen, and withdrew to another part of the island.[34] The settlement of San José de Oruña was later established by Antonio de Berrío on this land in 1592.[16][36] Shortly thereafter the English sailor Sir Walter Raleigh arrived in Trinidad on 22 March 1595 in search of the long-rumoured "El Dorado" ('City of Gold') supposedly located in South America.[36] He attacked San José, captured and interrogated Antonio de Berrío, and obtained much information from him and from the Cacique Topiawari; Raleigh then went on his way, and Spanish authority was restored.[40]
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+ Meanwhile, there were numerous attempts by European powers to settle Tobago during the 1620-40s, with the Dutch, English and Courlanders (people from the Duchy of Courland and Semigallia, now part of Latvia) all attempting to colonise the island with little success.[41][42] From 1654 the Dutch and Courlanders managed to gain a more secure foothold, later joined by several hundred French settlers.[41] A plantation economy developed based on the production of sugar, indigo and rum, worked by large numbers of African slaves who soon came to vastly outnumber the European colonists.[42][41] Large numbers of forts were constructed as Tobago became a source of contention between France, Holland and Britain, with the island changing hands some 31 times prior to 1814, a situation exacerbated by widespread piracy.[42] The British managed to hold Tobago from 1762–1781, whereupon it was captured by the French, who ruled until 1793 when Britain re-captured the island.[42]
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+ The 17th century on Trinidad passed largely without major incident, but sustained attempts by the Spaniards to control and rule over the Amerindians were often fiercely resisted.[36] In 1687 the Catholic Catalan Capuchin friars were given responsibility for the conversions of the indigenous people of Trinidad and the Guianas.[36] They founded several missions in Trinidad, supported and richly funded by the state, which also granted encomienda right to them over the native peoples, in which the native peoples were forced to provide labour for the Spanish.[36] One such mission was Santa Rosa de Arima, established in 1789, when Amerindians from the former encomiendas of Tacarigua and Arauca (Arouca) were relocated further west.[citation needed] Escalating tensions between the Spaniards and Amerindians culminated in violence 1689, when Amerindians in the San Rafael encomienda rebelled and killed several priests, attacked a church, and killed the Spanish governor José de León y Echales. Among those killed in the governor's party was Juan Mazien de Sotomayor, missionary priest to the Nepuyo villages of Cuara, Tacarigua and Arauca.[citation needed] The Spanish retaliated severely, slaughtering hundreds of native peoples in an event that became known as the Arena massacre.[36] As a result of this, continuing Spanish slave-raiding, and the devastating impact of introduced disease to which they had no immunity, the native population was virtually wiped out by the end of the following century.[43][36]
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+ During this period Trinidad was an island province belonging to the Viceroyalty of New Spain, together with Central America, present-day Mexico and the southwestern United States.[44] In 1757 the capital was moved from San José de Oruña to Puerto de España (modern Port of Spain) following several pirate attacks.[45] However the Spanish never made any concerted effort to colonise the islands; Trinidad in this period was still mostly forest, populated by a few Spaniards with a handful of slaves and a few thousand Amerindians.[44] Indeed, the population in 1777 was only 1,400, and Spanish colonisation in Trinidad remained tenuous.[citation needed]
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+
25
+ Since Trinidad was considered underpopulated, Roume de St. Laurent, a Frenchman living in Grenada, was able to obtain a Cédula de Población from the Spanish king Charles III on 4 November 1783.[46] A Cédula de Población had previously been granted in 1776 by the king, but had not shown results, and therefore the new Cédula was more generous.[16] It granted free land and tax exemption for 10 years to Roman Catholic foreign settlers who were willing to swear allegiance to the King of Spain.[16] The Spanish also gave many incentives to lure settlers to the island, including exemption from taxes for ten years and land grants in accordance with the terms set out in the Cédula.[47] The land grant was 30 fanegas (13 hectares/32 acres) for each free man, woman and child and half of that for each slave that they brought with them. The Spanish sent a new governor, José María Chacón, to implement the terms of the new cédula.[46]
26
+
27
+ It was fortuitous that the Cédula was issued only a few years before the French Revolution. During that period of upheaval, French planters with their slaves, free coloureds and mulattos from the neighbouring islands of Martinique, Saint Lucia, Grenada, Guadeloupe and Dominica migrated to Trinidad, where they established an agriculture-based economy (sugar and cocoa).[44] These new immigrants established local communities in Blanchisseuse, Champs Fleurs, Paramin,[48] Cascade, Carenage and Laventille.
28
+
29
+ As a result, Trinidad's population jumped to over 15,000 by the end of 1789, and by 1797 the population of Port of Spain had increased from under 3,000 to 10,422 in just five years, with a varied population of mixed race individuals, Spaniards, Africans, French republican soldiers, retired pirates and French nobility.[44] The total population of Trinidad was 17,718, of which 2,151 were of European ancestry, 4,476 were "free blacks and people of colour", 10,009 were enslaved people and 1,082 Amerindians.[citation needed] The sparse settlement and slow rate of population-increase during Spanish rule (and even later during British rule) made Trinidad one of the less populated colonies of the West Indies, with the least developed plantation infrastructure.[49]
30
+
31
+ The British had begun to take a keen interest in Trinidad, and in 1797 a British force led by General Sir Ralph Abercromby launched an invasion of Trinidad.[16][50] His squadron sailed through the Bocas and anchored off the coast of Chaguaramas. Seriously outnumbered, Governor Chacón decided to capitulate to British without fighting.[50] Trinidad thus became a British crown colony, with a largely French-speaking population and Spanish laws.[44] British rule was later formalised under the Treaty of Amiens (1802).[16][50] The colony's first British governor was Thomas Picton, however his heavy-handed approach to enforcing British authority, including the use of torture and arbitrary arrest, led to his being recalled.[50]
32
+
33
+ British rule led to an influx of settlers from the United Kingdom and the British colonies of the Eastern Caribbean. English, Scots, Irish, German and Italian families arrived, as well as some free blacks known as 'Merikins' who had fought for Britain in the War of 1812 and were granted land in southern Trinidad.[51][52][53] Under British rule, new states were created and the importation of slaves increased, however by this time support for abolitionism had vastly increased and in England the slave trade was under attack.[49][54] Slavery was abolished in 1833, after which former slaves served an "apprenticeship" period. In 1837 Daaga, a West African slave trader who had been captured by Portuguese slavers and later rescued by the British navy, was conscripted into the local regiment. Daaga and a group of his compatriots mutinied at the barracks in St Joseph and set out eastward in an attempt to return to their homeland. The mutineers were ambushed by a militia unit just outside the town of Arima. The revolt was crushed at the cost of some 40 dead, and Daaga and his party were later executed at St Joseph.[55] The apprenticeship system ended on 1 August 1838 with full emancipation.[16][53] An overview of the populations statistics in 1838, however, clearly reveals the contrast between Trinidad and its neighbouring islands: upon emancipation of the slaves in 1838, Trinidad had only 17,439 slaves, with 80% of slave owners having enslaved fewer than 10 people each.[56] In contrast, at twice the size of Trinidad, Jamaica had roughly 360,000 slaves.[57]
34
+
35
+ After the African slaves were emancipated many refused to continue working on the plantations, often moving out to urban areas such as Laventille and Belmont to the east of Port of Spain.[53] As a result, a severe agricultural labour shortage emerged; the British filled this gap by instituting a system of indentureship. Various nationalities were contracted under this system, including Indians, Chinese, and Portuguese.[58] Of these, the East Indians were imported in the largest numbers, starting from 1 May 1845, when 225 Indians were brought in the first shipment to Trinidad on the Fatel Razack, a Muslim-owned vessel.[53][59] Indentureship of the Indians lasted from 1845 to 1917, during which time more than 147,000 Indians came to Trinidad to work on sugarcane plantations.[16][60]
36
+
37
+ Indentureship contracts were sometimes exploitative, to such an extent that historians such as Hugh Tinker were to call it "a new system of slavery". Despite these descriptions, it was not truly a new form of slavery, as workers were paid, contracts were finite, and the idea of an individual being another's property had been eliminated when slavery was abolished.[61] In addition, employers of indentured labour had no legal right to flog or whip their workers; the main legal sanction for the enforcement of the indenture laws was prosecution in the courts, followed by fines or (more likely) jail sentences.[62] People were contracted for a period of five years, with a daily wage as low as 25 cents in the early 20th century, and they were guaranteed return passage to India at the end of their contract period. However, coercive means were often used to retain labourers, and the indentureship contracts were soon extended to 10 years from 1854 after the planters complained that they were losing their labour too early.[49][53] In lieu of the return passage, the British authorities soon began offering portions of land to encourage settlement, and by 1902, more than half of the sugar cane in Trinidad was being produced by independent cane farmers; the majority of which were Indians.[63] Despite the trying conditions experienced under the indenture system, about 90% of the Indian immigrants chose, at the end of their contracted periods of indenture, to make Trinidad their permanent home.[63] East Indians entering the colony were also subject to certain crown laws which segregated them from the rest of Trinidad's population, such as the requirement that they carry a pass with them if they left the plantations, and that if freed, they carry their "Free Papers" or certificate indicating completion of the indenture period.[64]
38
+
39
+ Few Indians settled on Tobago however, and the descendants of African slaves continued to form the majority of the island's population. An ongoing economic slump in the middle-to-late 19th century caused
40
+ widespread poverty.[65] Discontent erupted into rioting on the Roxborough plantation in 1876, in an event known as the Belmanna Uprising after a policeman who was killed.[65] The British eventually managed to restore control, however as a result of the disturbances Tobago's Legislative Assembly voted to dissolve itself and the island became a Crown colony in 1877.[65] With the sugar industry in a state of near-collapse and the island no longer profitable, the British attached Tobago to their Trinidad colony in 1899.[16][66][67]
41
+
42
+ In 1903, a protest against the introduction of new water rates in Port of Spain erupted into rioting; 18 people were shot dead, and the Red House (the government headquarters) was damaged by fire.[66] A local elected assembly with some limited powers was introduced in 1913.[66] Economically Trinidad and Tobago remained a predominantly agricultural colony; alongside sugarcane, the cacao (cocoa) crop also contributed greatly to economic earnings in the late 19th and early 20th centuries.
43
+
44
+ In November 1919, the dockworkers went on strike over bad management practises, low wages compared to a higher cost of living.[68] Strikebreakers were brought in to keep a minimum of goods moving through the ports. On December 1, 1919, the striking dockworkers rushed the harbour and chased off the strikebreakers.[68] They then proceeded to march on the government buildings in Port of Spain. Other unions and workers, many with the same grievances, joined the dock worker's strike making it a General Strike.[68] Violence broke out and was only put down with help from the sailors of British Naval ship HMS Calcutta. The unity brought upon by the strike was the first time of cooperation between the various ethnic groups of the time.[69] Historian Brinsley Samaroo says that the 1919 strikes "seem to indicate that there was a growing class consciousness after the war and this transcended racial feelings at times."[69]
45
+
46
+ However, in the 1920s, the collapse of the sugarcane industry, concomitant with the failure of the cocoa industry, resulted in widespread depression among the rural and agricultural workers in Trinidad, and encouraged the rise of a labour movement. Conditions on the islands worsened in the 1930s with the onset of the Great Depression, with an outbreak of labour riots occurring in 1937 which resulted in several deaths.[70] The labour movement aimed to unite the urban working class and agricultural labour class; the key figures being Arthur Cipriani, who led the Trinidad Workingmen's Association (TWA), and Tubal Uriah "Buzz" Butler of the British Empire Citizens' and Workers' Home Rule Party.[70] As the movement developed calls for greater autonomy from British colonial rule became widespread; this effort was severely undermined by the British Home Office and by the British-educated Trinidadian elite, many of whom were descended from the plantocracy class.
47
+
48
+ Petroleum had been discovered in 1857, but became economically significant only in the 1930s and afterwards as a result of the collapse of sugarcane and cocoa, and increasing industrialisation.[71]
49
+ [72][73] By the 1950s petroleum had become a staple in Trinidad's export market, and was responsible for a growing middle class among all sections of the Trinidad population. The collapse of Trinidad's major agricultural commodities, followed by the Depression, and the rise of the oil economy, led to major changes in the country's social structure.
50
+
51
+ The presence of American military bases in Chaguaramas and Cumuto in Trinidad during World War II had a profound effect on society. The Americans vastly improved the infrastructure on Trinidad and provided many locals with well-paying jobs; however the social effects of having so many young soldiers stationed on the island, as well as their often unconcealed racial prejudice, caused resentment.[66] The Americans left in 1961.[74]
52
+
53
+ In the post-war period the British began a process of decolonisation across the British Empire. In 1945 universal suffrage was introduced to Trinidad and Tobago.[16][66] Political parties emerged on the island, however these were largely divided along racial lines: Afro-Trinidadians and Tobagonians primarily supported the People's National Movement (PNM), formed in 1956 by Eric Williams, with Indo-Trinidadians and Tobagonians mostly supporting the People's Democratic Party (PDP), formed in 1953 by Bhadase Sagan Maraj,[75] which later merged into the Democratic Labour Party (DLP) in 1957.[76] Britain's Caribbean colonies formed the West Indies Federation in 1958 as a vehicle for independence, however the Federation dissolved after Jamaica withdrew following a membership referendum in 1961. The government of Trinidad and Tobago subsequently chose to seek independence from the United Kingdom on its own.[77]
54
+
55
+ Trinidad and Tobago gained its independence from the United Kingdom on 31 August 1962.[16][73] Elizabeth II remained head of state as Queen of Trinidad and Tobago, represented locally by Governor-General Solomon Hochoy. Eric Williams of the PNM, a noted historian and intellectual widely regarded as The Father of The Nation, became the first Prime Minister, serving in that capacity uninterrupted until 1981.[16] The dominant figure in the opposition in the early independence years was Rudranath Capildeo of the DLP. The 1960s saw the rise of a Black Power movement, inspired in part by the civil rights movement in the United States. Protests and strikes became common, with events coming to head in April 1970 when police shot dead a protester named Basil Davis.[76] Fearing a breakdown of law and order, Prime Minister Williams declared a state of emergency and arrested many of the Black Power leaders. Some army leaders who were sympathetic to the Black Power movement, notably Raffique Shah and Rex Lassalle, attempted to mutiny; however, this was quashed by the Trinidad and Tobago Coast Guard.[76] Williams and the PNM retained power, largely due to divisions in the opposition.[76]
56
+
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+ In 1963 Tobago was struck by Hurricane Flora, which killed 30 people and resulted in enormous destruction across the island.[78] Partly as a result of this, tourism came to replace agriculture as the island's main income earner in the subsequent decades.[78]
58
+
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+ Between the years 1972 and 1983, the country profited greatly from the rising price of oil and the discovery of vast new oil deposits in its territorial waters, resulting in an economic boom that increased living standards greatly.[16][76] In 1976 the country became a republic within the Commonwealth, though it retained the Judicial Committee of the Privy Council as its final appellate court.[16] The position of governor-general was replaced with that of President; Ellis Clarke was the first to hold this largely ceremonial role.[79] Tobago was granted limited self-rule with the creation of the Tobago House of Assembly in 1980.[65]
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+
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+ Williams died in 1981, being replaced by George Chambers who led the country until 1986. By this time a fall in the price of oil had resulted in a recession, causing rising inflation and unemployment.[80] The main opposition parties united under the banner of National Alliance for Reconstruction (NAR) and won the 1986 Trinidad and Tobago general election, with NAR leader A. N. R. Robinson becoming the new Prime Minister.[81][76] Robinson was unable to hold together the fragile NAR coalition, and social unrest was caused by his economic reforms, such as devaluing the currency and implementing an International Monetary Fund Structural Adjustment Program.[16] In 1990 114 members of the Jamaat al Muslimeen, led by Yasin Abu Bakr (formerly known as Lennox Phillip) stormed the Red House (the seat of Parliament), and Trinidad and Tobago Television, the only television station in the country at the time, holding Robinson and country's government hostage for six days before surrendering.[82] The coup leaders were promised amnesty, but upon their surrender they were then arrested, but later released after protracted legal wrangling.[58]
62
+
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+ The PNM under Patrick Manning returned to power following the 1991 Trinidad and Tobago general election.[16] Hoping to capitalise on an improvement in the economy, Manning called an early election in 1995, however, this resulted in a hung parliament. Two NAR representatives backed the opposition United National Congress (UNC), which had split off from the NAR in 1989, and they thus took power under Basdeo Panday, who became the country's first Indo-Trinidadian Prime Minister.[16][80][83] After a period of political confusion caused by a series of inconclusive election results, Patrick Manning returned to power in 2001, retaining that position until 2010.[16]
64
+
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+ Since 2003 the country entered a second oil boom, and petroleum, petrochemicals and natural gas continue to be the backbone of the economy. Tourism and the public service are the mainstay of the economy of Tobago, though authorities have attempted to diversify the island's economy.[84] A corruption scandal resulted in Manning's defeat by the newly formed People's Partnership coalition in 2010, with Kamla Persad-Bissessar becoming the country's first female Prime Minister.[85][86][87] However, corruption allegations bedevilled the new administration, and the PP were defeated in 2015 by the PNM under Keith Rowley.[88][89]
66
+
67
+ Trinidad and Tobago is situated between 10° 2' and 11° 12' N latitude and 60° 30' and 61° 56' W longitude, with the Caribbean Sea to the north, the Atlantic Ocean to the east and south, and the Gulf of Paria to the west. It is located in the far south-east of the Caribbean region, with the island of Trinidad being just 11 kilometres (6.8 mi) off the coast of Venezuela in mainland South America across the Columbus Channel.[16] Covering an area of 5,128 km2 (1,980 sq mi),[90] the country consists of two main islands, Trinidad and Tobago, separated by a 20m (30 km) strait, plus a number of much smaller islands, including Chacachacare, Monos, Huevos, Gaspar Grande (or Gasparee), Little Tobago, and Saint Giles Island.[16]
68
+
69
+ Trinidad is 4,768 km2 (1,841 sq mi) in area (comprising 93.0% of the country's total area) with an average length of 80 kilometres (50 mi) and an average width of 59 kilometres (37 mi). Tobago has an area of about 300 km2 (120 sq mi), or 5.8% of the country's area, is 41 km (25 mi) long and 12 km (7.5 mi) at its greatest width. Trinidad and Tobago lie on the continental shelf of South America, and are thus geologically considered to lie entirely in South America.[16]
70
+
71
+ The terrain of the islands is a mixture of mountains and plains.[15] On Trinidad the Northern Range runs parallel with the north coast, and contains the country's highest peak (El Cerro del Aripo), which is 940 metres (3,080 ft) above sea level[15] and second highest (El Tucuche, 936 metres (3,071 ft)).[16] The rest of the island is generally flatter, excluding the Central Range and Montserrat Hills in the centre of the island and the Southern Range and Trinity Hills in the south. The east coast is noted for its beaches, most notably Manzanilla Beach. The island contains several large swamp areas, such as the Caroni Swamp and the Nariva Swamp.[16] Major bodies of water on Trinidad include the Hollis Reservoir, Navet Reservoir, Caroni Reservoir. Trinidad is made up of a variety of soil types, the majority being fine sands and heavy clays. The alluvial valleys of the Northern Range and the soils of the East–West Corridor are the most fertile.[91][citation needed] Trinidad is also notable for containing Pitch Lake, the largest natural reservoir of asphalt in the world.[15][16] Tobago contains a flat plain in its south-west, with the eastern half of the island being more mountainous, culminating in Pigeon Peak, the island's highest point at 550 metres (1,800 ft).[92] Tobago also contains several coral reefs off its coast.[16]
72
+
73
+ The majority of the population reside on the island of Trinidad, and this is thus the location of largest towns and cities. There are four major municipalities in Trinidad: the capital Port of Spain, San Fernando, Arima and Chaguanas. The main town on Tobago is Scarborough.
74
+
75
+ The Northern Range consists mainly of Upper Jurassic and Cretaceous metamorphic rocks. The Northern Lowlands (the East–West Corridor and Caroni Plain) consist of younger shallow marine clastic sediments. South of this, the Central Range fold and thrust belt consists of Cretaceous and Eocene sedimentary rocks, with Miocene formations along the southern and eastern flanks. The Naparima Plain and the Nariva Swamp form the southern shoulder of this uplift.[citation needed]
76
+
77
+ The Southern Lowlands consist of Miocene and Pliocene sands, clays, and gravels. These overlie oil and natural gas deposits, especially north of the Los Bajos Fault. The Southern Range forms the third anticlinal uplift. The rocks consist of sandstones, shales, siltstones and clays formed in the Miocene and uplifted in the Pleistocene. Oil sands and mud volcanoes are especially common in this area.[citation needed]
78
+
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+ Trinidad and Tobago has a maritime tropical climate.[15][16] There are two seasons annually: the dry season for the first five months of the year, and the rainy season in the remaining seven of the year. Winds are predominantly from the northeast and are dominated by the northeast trade winds. Unlike many Caribbean islands Trinidad and Tobago lies outside the main hurricane alleys; nevertheless, the island of Tobago was struck by Hurricane Flora on September 30, 1963. In the Northern Range of Trinidad, the climate is often cooler than that of the sweltering heat of the plains below, due to constant cloud and mist cover, and heavy rains in the mountains.
80
+
81
+ Record temperatures for Trinidad and Tobago are 39 °C (102 °F)[93] for the high in Port of Spain, and a low of 12 °C (54 °F).[94]
82
+
83
+ Because Trinidad and Tobago lies on the continental shelf of South America, and in ancient times were physically connected to the South American mainland, its biological diversity is unlike that of most other Caribbean islands, and has much more in common with that of Venezuela.[95] The main ecosystems are: coastal and marine (coral reefs, mangrove swamps, open ocean and seagrass beds); forest; freshwater (rivers and streams); karst; man-made ecosystems (agricultural land, freshwater dams, secondary forest); and savanna. On 1 August 1996, Trinidad and Tobago ratified the 1992 Rio Convention on Biological Diversity, and it has produced a biodiversity action plan and four reports describing the country's contribution to biodiversity conservation. These reports formally acknowledged the importance of biodiversity to the well-being of the country's people through provision of ecosystem services.[96]
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+
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+ Information about vertebrates is good, with 472 bird species (2 endemics), about 100 mammals, about 90 reptiles (a few endemics), about 30 amphibians (including several endemics), 50 freshwater fish and at least 950 marine fish.[97] Notable mammal species include the ocelot, manatee, collared peccary (known as the quenk locally), agouti, lappe, red brocket deer, otter, weeper capuchin and red howler monkey; there are also some 70 species of bat, including the vampire bat and fringe-lipped bat.[16][98] Amongst the reptiles, the spectacled caiman is the largest, sometimes growing up to 3m.[95] There are also 47 species of snake, including four venomous species, lizards such as the gecko, iguana, matte lizard and also several species of turtle.[16][99] are present. Of the amphibians, the golden tree frog is endemic to Trinidad.[99] Marine life is abundant, with several species of sea urchin, coral, lobster, anemone, starfish, manta ray, dolphin, porpoise and whale shark present in the islands' waters.[100] The introduced lionfish is viewed as a pest, as it eats many native species of fish and has no natural predators; efforts are currently underway to cull the numbers of this species.[100]
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+
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+ Trinidad and Tobago is noted particularly for its large number of bird species, and is a popular destination for bird watchers. Notable species include the scarlet ibis, cocrico, egret, shiny cowbird, bananaquit, oilbird and various species of honeycreeper, trogon, toucan, parrot, tanager, woodpecker, antbird, kites, hawks, boobies, pelicans and vultures; there are also 17 species of hummingbird, including the tufted coquette which is the world's third smallest.[101]
88
+
89
+ Information about invertebrates is dispersed and very incomplete. About 650 butterflies,[97] at least 672 beetles (from Tobago alone)[102] and 40 corals[97] have been recorded.[97] Other notable invertebrates include the cockroach, leaf-cutter ant and numerous species of mosquitoes, termites, spiders and tarantulas.
90
+
91
+ Although the list is far from complete, 1,647 species of fungi, including lichens, have been recorded.[103][104][105] The true total number of fungi is likely to be far higher, given the generally accepted estimate that only about 7% of all fungi worldwide have so far been discovered.[106] A first effort to estimate the number of endemic fungi tentatively listed 407 species.[107]
92
+
93
+ Information about micro-organisms is dispersed and very incomplete. Nearly 200 species of marine algae have been recorded.[97] The true total number of micro-organism species must be much higher.
94
+
95
+ Thanks to a recently published checklist, plant diversity in Trinidad and Tobago is well documented with about 3,300 species (59 endemic) recorded.[97] Despite significant felling, forests still cover about 40% of the country, and there are about 350 different species of tree.[95] A notable tree is the manchineel which is extremely poisonous to humans, and even just touching its sap can cause severe blistering of the skin; the tree is often covered with warning signs.
96
+
97
+ Trinidad and Tobago is a republic with a two-party system and a bicameral parliamentary system based on the Westminster System.[15]
98
+
99
+ The head of state of Trinidad and Tobago is the President, currently Paula Mae Weekes.[15] This largely ceremonial role replaced that of the Governor-General (representing the Monarch of Trinidad and Tobago) upon Trinidad and Tobago's becoming a republic in 1976.[16] The head of government is the Prime Minister, currently Keith Rowley.[15] The President is elected by an Electoral college consisting of the full membership of both houses of Parliament.
100
+ The Prime Minister is elected following a general election which takes place every five years. The President is required to appoint the leader of the party who in his or her opinion has the most support of the members of the House of Representatives to this post; this has generally been the leader of the party which won the most seats in the previous election (except in the case of the 2001 General Elections).[16]
101
+
102
+ Since 1980 Tobago has also had its own elections, separate from the general elections. In these elections, members are elected and serve in the unicameral Tobago House of Assembly.[108][15][16]
103
+
104
+ Parliament consists of the Senate (31 seats) and the House of Representatives (41 seats, plus the Speaker).[109][15] The members of the Senate are appointed by the president; 16 Government Senators are appointed on the advice of the Prime Minister, six Opposition Senators are appointed on the advice of the Leader of the Opposition, currently Kamla Persad-Bissessar, and nine Independent Senators are appointed by the President to represent other sectors of civil society. The 41 members of the House of Representatives are elected by the people for a maximum term of five years in a "first past the post" system.
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+
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+ Trinidad is split into 14 regional corporations and municipalities, consisting of nine regions and five municipalities, which have a limited level of autonomy.[15][16] The various councils are made up of a mixture of elected and appointed members. Elections are held every three years.[citation needed] The country was formerly divided into counties.
107
+
108
+ The two main parties are the People's National Movement (PNM) and the United National Congress (UNC); another recent party was the Congress of the People (COP). Support for these parties appears to fall along ethnic lines, with the PNM consistently obtaining a majority of Afro-Trinidadian vote, and the UNC gaining a majority of Indo-Trinidadian support.
109
+
110
+ The Trinidad and Tobago Defence Force (TTDF) is the military organisation responsible for the defence of the twin island Republic of Trinidad and Tobago.[15] It consists of the Regiment, the Coast Guard, the Air Guard and the Defence Force Reserves. Established in 1962 after Trinidad and Tobago's independence from the United Kingdom, the TTDF is one of the largest military forces in the Anglophone Caribbean.[citation needed]
111
+
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+ Its mission statement is to "defend the sovereign good of The Republic of Trinidad and Tobago, contribute to the development of the national community and support the State in the fulfillment of its national and international objectives". The Defence Force has been engaged in domestic incidents, such as the 1990 Coup Attempt, and international missions, such as the United Nations Mission in Haiti between 1993 and 1996.
113
+
114
+ In 2019, Trinidad and Tobago signed the UN treaty on the Prohibition of Nuclear Weapons.[110]
115
+
116
+ Trinidad and Tobago maintains close relations with its Caribbean neighbours and major North American and European trading partners. As the most industrialised and second-largest country in the Anglophone Caribbean, Trinidad and Tobago has taken a leading role in the Caribbean Community (CARICOM), and strongly supports CARICOM economic integration efforts. It also is active in the Summit of the Americas process and supports the establishment of the Free Trade Area of the Americas, lobbying other nations for seating the Secretariat in Port of Spain.[citation needed]
117
+
118
+ As a member of CARICOM, Trinidad and Tobago strongly backed efforts by the United States to bring political stability to Haiti, contributing personnel to the Multinational Force in 1994. After its 1962 independence, Trinidad and Tobago joined the United Nations and Commonwealth of Nations. In 1967 it became the first Commonwealth country to join the Organization of American States (OAS).[111] In 1995 Trinidad played host to the inaugural meeting of the Association of Caribbean States and has become the seat of this 35-member grouping, which seeks to further economic progress and integration among its states. In international forums, Trinidad and Tobago has defined itself as having an independent voting record, but often supports US and EU positions.[citation needed]
119
+
120
+ Trinidad and Tobago has in recent decades suffered from a relatively high crime rate;[112][113] there are currently roughly 500 murders per year.[114][76] The country is a noted transshipment centre for the trafficking of illegal drugs from South America to the rest of the Caribbean and beyond to North America.[115] Some estimates put the size of the 'hidden economy' as high as 20–30% of measured GDP.[116]
121
+
122
+ Though there have been no terrorism-related incidents in the country since the 1990 Islamic coup attempt, Trinidad and Tobago remains a potential target; for example, in February 2018 a plan to attack the Carnival was foiled by police.[113] It is estimated that roughly 100 citizens of the country have traveled to the Middle East to fight for Islamic State.[112][113] In 2017 the government adopted a counter-terrorism and extremism strategy.[113]
123
+
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+ The country's prison administration is the Trinidad and Tobago Prison Service (TTPS), it is under the control of the Commissioner of Prisons Gerard Wilson, located in Port-of-Spain.[117] The prison population rate is 292 people per 100,000. The total prison population, including pre-trial detainees and remand prisoners, is 3,999 prisoners. The population rate of pre-trial detainees and remand prisoners is 174 per 100,000 of the national population (59.7% of the prison population). In 2018, the female prison population rate is 8.5 per 100,000 of the national population (2.9% of the prison population). Prisoners that are minors makes up 1.9% of the prison population and foreigners prisoners make 0.8% of the prison population. The occupancy level of Trinidad and Tobago's prison system is at 81.8% capacity.[117] Trinidad and Tobago has nine prison establishments; Golden Grove Prison, Maximum Security Prison, Port of Spain Prison, Eastern Correctional Rehabilitation Centre, Remand Prison, Tobago Convict Prison, Carrera Convict Island Prison, Women's Prison and Youth Training and Rehabilitation Centre.[118] Trinidad and Tobago also use labor yards as prisons, or means of punishment.[119]
125
+
126
+ The population of the country currently stands at 1,363,985 (July 2019 est.).[citation needed]
127
+
128
+ The ethnic composition of Trinidad and Tobago reflects a history of conquest and immigration.[120] While the earliest inhabitants were of Amerindian heritage, the two dominant groups in the country are now those of South Asian and of African heritage. Indo-Trinidadian and Tobagonians make up the country's largest ethnic group (approximately 35.4%);[15] they are primarily the descendants of indentured workers from South Asia (mostly from India), brought to replace freed African slaves who refused to continue working on the sugar plantations. Through cultural preservation many residents of Indian descent continue to maintain traditions from their ancestral homeland. Indo-Trinidadians reside primarily on Trinidad; as of the 2011 census only 2.5% of Tobago's population was of Indian descent.[121]
129
+
130
+ Afro-Trinidadians and Tobagonians make up the country's second largest ethnic group, with approximately 34.2% of the population identifying as being of African descent.[15] The majority of people of an African background are the descendants of slaves forcibly transported to the islands from as early as the 16th century. This group constitute the majority on Tobago, at 85.2%.[121]
131
+
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+ The bulk of the rest of the population are those who identify as being of mixed heritage.[15] There are also small but significant minorities of people of Amerindian, European, Chinese, and Arab descent. Arima on Trinidad is a noted centre of Amerindian culture.[16]
133
+
134
+ English is the country's official language (the local variety of standard English is Trinidadian and Tobagonian English or more properly, Trinidad and Tobago Standard English, abbreviated as "TTSE"), but the main spoken language is either of two English-based creole languages (Trinidadian Creole or Tobagonian Creole), which reflects the Amerindian, European, African, and Asian heritage of the nation. Both creoles contain elements from a variety of African languages; Trinidadian English Creole, however, is also influenced by French and French Creole (Patois).[122]
135
+
136
+ A majority of the early Indian immigrants spoke the Bhojpuri and Awadhi dialect of Hindustani (Hindi-Urdu), which later formed into Trinidadian Hindustani (Hindi-Urdu), which became the lingua franca of Indo-Trinidadian and Tobagonians. From 1935 Indian films began showing to audiences in Trinidad; most of these were in the Standard Hindustani (Hindi-Urdu) dialect and this modified Trinidadian Hindustani slightly by adding Standard Hindu and Urdu phrases and vocabulary to Trinidadian Hindustani. Indian films also revitalised Hindustani among Indo-Trinidadian and Tobagonians.[123] Around the mid to late 1970s the lingua franca of Indo-Trinidaian and Tobagonians switched from Trinidadian Hindustani to a sort of Hindinised version of English. Today Hindustani survives on through Indo-Trinidadian and Tobagonian musical forms such as, Bhajan, Indian classical music, Indian folk music, Filmi, Pichakaree, Chutney, Chutney soca, and Chutney parang. Presently there are about 26,000 people, which is 5.53% of the Indo-Trinidadian and Tobagonian population, who speak Trinidadian Hindustani. Many Indo-Trinidadians and Tobagonians today speak a type of Hinglish that consist of Trinidadian and Tobagonian English that is heavily laced with Trinidadian Hindustani vocabulary and phrases and many Indo-Trinidadians and Tobagonians can recite phrases or prayers in Hindustani today. There are many places in Trinidad and Tobago that have names of Hindustani origin. Some phrases and vocabulary have even made their way into the mainstream English and English Creole dialects of the country.[124][125][126][127][128]
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+
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+ The Chinese language first came to Trinidad and Tobago in 1806, when the British had brought Chinese labourers in order to determine if they were fit to use as labourers after the abolition of slavery.[citation needed] About 2,645 Chinese immigrants arrived in Trinidad as indentured labour between 1853 and 1866.[citation needed] A majority of the people who immigrated in the 19th century were from southern China and spoke the Hakka and Yue dialects of Chinese. In the 20th century after the years of indentureship up to the present-day more Chinese people have immigrated to Trinidad and Tobago for business and they speak the dialects of the indenturees along with other Chinese dialects, such as Mandarin and Min.[125][129] J. Dyer Ball, writing in 1906, says: "In Trinidad there were, about twenty years ago, 4,000 or 5,000 Chinese, but they have decreased to probably about 2,000 or 3,000, [2,200 in 1900]. They used to work in sugar plantations, but are now principally shopkeepers, as well as general merchants, miners and railway builders,
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+ etc."[130]
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+
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+ The indigenous languages were Yao on Trinidad and Karina on Tobago, both Cariban, and Shebaya on Trinidad, which was Arawakan.[125]
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+
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+ According to the 2011 census,[3] Roman Catholics were the largest single religious group in Trinidad and Tobago with 21.60% of the total population. The Pentecostal/Evangelical/Full Gospel denominations were the third largest group with 12.02% of the population. The remaining population is made of various Christian denominations (Spiritual Shouter Baptists (5.67%), Anglicans (5.67%), Seventh-day Adventists (4.09%), Presbyterians or Congregationalists (2.49%), Jehovah's Witnesses (1.47%), other Baptists (1.21%), Methodists (0.65%) and the Moravian Church (0.27%)). Respondents who did not state a religious affiliation represented 11.1% of the population, with 2.18% declaring themselves Irreligious.
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+ Hindus were the second largest group with 18.15%.[3] Hinduism is practiced throughout the country and Diwali is a public holiday, and other Hindu holidays are also widely celebrated.
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+
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+ Muslims represent 4.97% of the population.[3] Eid al-Fitr is a public holiday and Eid al-Adha, Mawlid, Hosay, and other Muslim holidays are also celebrated. There has also been a Jewish community on the islands for many centuries, however their numbers have never been large, with a 2007 estimating putting the Jewish population at 55 individuals.[131][132]
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+ African-derived or Afrocentric religions are also practised, notably Trinidad Orisha (Yoruba) believers (0.9%) and Rastafarians (0.27%).[3] Various aspects of traditional obeah beliefs are still commonly practised on the islands.[50]
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+ Two African syncretic faiths, the Shouter or Spiritual Baptists and the Orisha faith (formerly called Shangos, a less than complimentary term)[citation needed] are among the fastest growing religious groups. Similarly, there is a noticeable increase in numbers of Evangelical Protestant and Fundamentalist churches usually lumped as "Pentecostal" by most Trinidadians, although this designation is often inaccurate. Sikhism, Jainism, Bahá'í, and Buddhism are practised by a minority of Indo-Trinidadian and Tobagonians. Several eastern religions such as Buddhism and Chinese folk religions such as Taoism and Confucianism are followed by Chinese Trinidadian and Tobagonian.
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+
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+ Children generally start pre-school at two and a half years but this is not mandatory. They are however, expected to have basic reading and writing skills when they commence primary school. Students proceed to a primary school at the age of five years. Seven years are spent in primary school. The seven classes of primary school consists of First Year and Second Year, followed by Standard One through Standard Five. During the final year of primary school, students prepare for and sit the Secondary Entrance Assessment (SEA) which determines the secondary school the child will attend.
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+
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+ Students attend secondary school for a minimum of five years, leading to the CSEC (Caribbean Secondary Education Certificate) examinations, which is the equivalent of the British GCSE O levels. Children with satisfactory grades may opt to continue high school for a further two-year period, leading to the Caribbean Advanced Proficiency Examinations (CAPE), the equivalent of GCE A levels. Both CSEC and CAPE examinations are held by the Caribbean Examinations Council (CXC). Public Primary and Secondary education is free for all, although private and religious schooling is available for a fee.
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+
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+ Tertiary education for tuition costs are provided for via GATE (The Government Assistance for Tuition Expenses), up to the level of the bachelor's degree, at the University of the West Indies (UWI), the University of Trinidad and Tobago (UTT), the University of the Southern Caribbean (USC), the College of Science, Technology and Applied Arts of Trinidad and Tobago (COSTAATT) and certain other local accredited institutions. Government also currently subsidises some Masters programmes. Both the Government and the private sector also provide financial assistance in the form of academic scholarships to gifted or needy students for study at local, regional or international universities.
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+ While women account for only 49% of the population, they constitute nearly 55% of the workforce in the country.[137]
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+ Trinidad and Tobago is the most developed nation and one of the wealthiest in the Caribbean and is listed in the top 40 (2010 information) of the 70 high-income countries in the world.[citation needed] Its gross national income per capita of US$20,070[138] (2014 gross national income at Atlas Method) is one of the highest in the Caribbean.[139] In November 2011, the OECD removed Trinidad and Tobago from its list of developing countries.[140] Trinidad's economy is strongly influenced by the petroleum industry. Tourism and manufacturing are also important to the local economy. Tourism is a growing sector, particular on Tobago, although proportionately it is much less important than in many other Caribbean islands. Agricultural products include citrus and cocoa. It also supplies manufactured goods, notably food, beverages, and cement, to the Caribbean region.
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+ Trinidad and Tobago is the leading Caribbean producer of oil and gas, and its economy is heavily dependent upon these resources.[16] Oil and gas account for about 40% of GDP and 80% of exports, but only 5% of employment.[15] Recent growth has been fuelled by investments in liquefied natural gas (LNG), petrochemicals, and steel. Additional petrochemical, aluminium, and plastics projects are in various stages of planning.
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+ The country is also a regional financial centre, and the economy has a growing trade surplus.[90] The expansion of Atlantic LNG over the past six years created the largest single-sustained phase of economic growth in Trinidad and Tobago. The nation is an exporter of LNG and supplied a total of 13.4 billion m3 in 2017. The largest markets for Trinidad and Tobago's LNG exports are Chile and the United States.[141]
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+ Trinidad and Tobago has transitioned from an oil-based economy to a natural gas based economy. In 2017, natural gas production totalled 18.5 billion m3, a decrease of 0.4% from 2016 with 18.6 billion m3 of production.[141] Oil production has decreased over the past decade from 7.1 million metric tonnes per year in 2007 to 4.4 million metric tonnes per year in 2017.[142] In December 2005, the Atlantic LNG's fourth production module or "train" for liquefied natural gas (LNG) began production. Train four has increased Atlantic LNG's overall output capacity by almost 50% and is the largest LNG train in the world at 5.2 million tons/year of LNG.[citation needed]
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+ Trinidad and Tobago is far less dependent on tourism than many other Caribbean countries and territories, with the bulk of tourist activity occurring on Tobago.[16] The government has made efforts to boost this sector in recent years.[16]
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+ Historically agricultural production (for example, sugar and coffee) dominated the economy, however this sector has been in steep decline since the 20th century and now forms just 0.4% of the country's GDP, employing 3.1% of the workforce.[15][16] Various fruits and vegetables are grown, such as cucumbers, eggplant, cassava, pumpkin, dasheen (taro) and coconut; fishing is still also commonly practised.[15]
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+ Trinidad and Tobago, in an effort to undergo economic transformation through diversification,[15] formed InvesTT in 2012 to serve as the country's sole investment promotion agency. This agency is aligned to the Ministry of Trade and Industry and is to be the key agent in growing the country's non-oil and gas sectors significantly and sustainably.[143]
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+ Trinidad and Tobago has a well developed communications sector. The telecommunications and broadcasting sectors generated an estimated TT$5.63 billion (US$0.88 billion) in 2014, which as a percentage of GDP equates to 3.1 percent. This represented a 1.9 percent increase in total revenues generated by this industry compared to last year. Of total telecommunications and broadcasting revenues, mobile voice services accounted for the majority of revenues with TT$2.20 billion (39.2 percent). This was followed by internet services which contributed TT$1.18 billion or 21.1 percent. The next highest revenue earners for the industry were fixed voice services and paid television services whose contributions totalled TT$0.76 billion and TT$0.70 billion respectively (13.4 percent and 12.4 percent). International voice services was next in line, generating TT$0.27 billion (4.7 percent) in revenues. Free-to Air radio and television services contributed TT$0.18 billion and TT$0.13 billion respectively (3.2 percent and 2.4 percent). Finally, other contributors included "other revenues" and "leased line services" with earnings of TT$0.16 billion and TT$0.05 billion respectively, with 2.8 percent and 0.9 percent.[144]
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+ There are several providers for each segment of the telecommunications market. Fixed Lines Telephone service is provided by Digicel, TSTT (operating as bmobile) and Cable & Wireless Communications operating as FLOW; cellular service is provided by TSTT (operating as bmobile) and Digicel whilst internet service is provided by TSTT, FLOW, Digicel, Green Dot and Lisa Communications.
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+ The Government of Trinidad and Tobago has recognised the creative industries as a pathway to economic growth and development. It is one of the newest, most dynamic sectors where creativity, knowledge and intangibles serve as the basic productive resource. In 2015, the Trinidad and Tobago Creative Industries Company Limited (CreativeTT) was established as a state agency under the Ministry of Trade and Industry with a mandate to stimulate and facilitate the business development and export activities of the Creative Industries in Trinidad and Tobago to generate national wealth, and, as such, the company is responsible for the strategic and business development of the three (3) niche areas and sub sectors currently under its purview – Music, Film and Fashion. MusicTT, FilmTT and FashionTT are the subsidiaries established to fulfil this mandate.
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+ The transport system in Trinidad and Tobago consists of a dense network of highways and roads across both major islands, ferries connecting Port of Spain with Scarborough and San Fernando, and international airports on both islands.[16] The Uriah Butler Highway, Churchill Roosevelt Highway and the Sir Solomon Hochoy Highway links the island of Trinidad together, whereas the Claude Noel Highway is the only major highway in Tobago. Public transportation options on land are public buses, private taxis and minibuses. By sea, the options are inter-island ferries and inter-city water taxis.[145]
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+ The island of Trinidad is served by Piarco International Airport located in Piarco, which opened on 8 January 1931.[citation needed] Elevated at 17.4 metres (57 ft) above sea level it comprises an area of 680 hectares (1,700 acres) and has a runway of 3,200 metres (10,500 ft). The airport consists of two terminals, the North Terminal and the South Terminal. The older South Terminal underwent renovations in 2009 for use as a VIP entrance point during the 5th Summit of the Americas. The North Terminal was completed in 2001, and consists of[146] 14-second-level aircraft gates with jetways for international flights, two ground-level domestic gates and 82 ticket counter positions.
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+ In 2008 the passenger throughput at Piarco International Airport was approximately 2.6 million. It is the seventh busiest airport in the Caribbean and the third busiest in the English-speaking Caribbean, after Sangster International Airport and Lynden Pindling International Airport.[citation needed] Caribbean Airlines, the national airline, operates its main hub at the Piarco International Airport and services the Caribbean, the United States, Canada and South America. The airline is wholly owned by the Government of Trinidad and Tobago. After an additional cash injection of US$50 million, the Trinidad and Tobago government acquired the Jamaican airline Air Jamaica on 1 May 2010, with a 6–12-month transition period to follow.[147]
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+ The Island of Tobago is served by the A.N.R. Robinson International Airport in Crown Point.[16] This airport has regular services to North America and Europe. There are regular flights between the two islands, with fares being heavily subsidised by the Government.
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+ Trinidad was formerly home to a railway network, however this was closed down in 1968.[148] There have been talks to build a new railway on the islands, though nothing yet has come of this.[149]
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+ The Strategic Plan for the Caribbean Community 2015–2019 was adopted by Trinidad and Tobago and the other members of the Caribbean Common Market (CARICOM) in 2014. The first of its kind, this document reflects a desire among countries to embrace a more profound regionalism, in order to reposition the Caribbean in an increasingly volatile global economy. The plan proposes mobilising funding from the public and private sectors to foster research and development (R&D) in the business sector. The plan outlines strategies for nurturing creativity, entrepreneurship, digital literacy and for making optimum use of available resources. It focuses on developing creative, manufacturing and service industries, with a special focus on tourism initially, natural resources and value-added products, agriculture and fisheries, to reduce dependence on food imports and foster sustainable fisheries, and energy efficiency.[150][151]
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+ Trinidad and Tobago is the region's leading exporter of oil and gas but imports of fossil fuels provided over 90% of the energy consumed by its CARICOM neighbours in 2008. This vulnerability led CARICOM to develop an Energy Policy which was approved in 2013. This policy is accompanied by the CARICOM Sustainable Energy Roadmap and Strategy (C-SERMS). Under the policy, renewable energy sources are to contribute 20% of the total electricity generation mix in member states by 2017, 28% by 2022 and 47% by 2027.[150]
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+ In 2014 Trinidad and Tobago was the third country in the world which emitted the most CO2 per capita after Qatar and Curacao according to the World Bank.[152] On average, each inhabitant produced 34.2 metric tons of CO2 in the atmosphere. In comparison, the world average was 5.0 tons per capita the same year.
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+ The Caribbean Community Climate Change Centre (CCCCC) has produced an implementation plan for 2011–2021 and carried out work to assess and build capacity in climate change mitigation and resilient development strategies. This work has been supported by the region's specialists, who have produced models for climate change and mitigation processes in Caribbean states. They also play a major advisor role to the divisions in ministries responsible for climate change. The growing frequency and intensity of hurricanes is of concern to all Caribbean nations. In 2012, Trinidad and Tobago had a 9% chance each year of being struck by a hurricane, according to estimates by the International Monetary Fund.[150][153]
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+ The two main bodies responsible for science, technology and innovation in Trinidad and Tobago are the Ministry of Science, Technology and Higher Education and the National Commission for Science and Technology.[150]
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+ In 2002, Trinidad and Tobago adopted Vision 2020. Like Jamaica's Vision 2030 (2009) and the Strategic Plan of Barbados for 2005–2025, Trinidad and Tobago's Vision 2020 accords central importance to harnessing science, technology and innovation (STI) to raise living standards and strengthen resilience to environmental shocks like hurricanes.[150]
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+ Trinidad and Tobago is one of the more affluent members of CARICOM, thanks to its natural resources. Despite this, it spent just 0.05% of GDP on R&D in 2012, according to the UNESCO Institute for Statistics. Even when the country was enjoying economic growth of 8% per annum in 2004, it devoted just 0.11% of GDP to R&D. Calculated in thousands of current Purchasing Power Parity (PPP) dollars, research expenditure actually dropped between 2009 and 2012 from 21 309 to 19 232. This corresponds to research expenditure of $PPP 65 per capita in 2009 and $PPP 45 in 2012.[150]
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+ Industrial R&D has declined since 2000, perhaps owing to the drop in research activity in the sugar sector. Whereas industrial R&D accounted for 24% of domestic research expenditure in 2004 and 29.5% in 2005, it had become almost non-existent by 2010.[150]
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+ The number of researchers in Trinidad and Tobago grew from 787 to 914 between 2009 and 2012. This corresponds in a rise from 595 to 683 in the number of researchers (head counts) per million inhabitants.[150]
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+ Scientific output grew between 2007 and 2011, according to Thomson Reuters' Web of Science (Science Citation Index Expanded) before contracting over the period 2012–2014. Trinidad and Tobago produced 109 publications per million population in 2014, behind Grenada (1,430), St Kitts and Nevis (730), Barbados (182) and Dominica (138) but ahead of the Bahamas (86), Belize (47) and Jamaica (42).[150]
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+ Between 2008 and 2014, scientists collaborated most with their peers from the United States (251 papers), United Kingdom (183), Canada (95), India (63) and Jamaica (43), according to the copublication record of Thomson Reuters. In turn, Jamaican scientists considered their counterparts from Trinidad and Tobago to be their fourth-closest collaborators (with 43 joint papers) after those from the United States, United Kingdom and Canada.[150]
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+ Between 2008 and 2013, Trinidad and Tobago registered 17 patents with the US Patent and Trademark Office (USPTO). This corresponds to 13% of the 134 patents registered by CARICOM members over this period. The top contributors were the Bahamas (34 patents) and Jamaica (22).[150]
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+ Trinidad and Tobago led CARICOM members for the value of high-tech exports in 2008 (US$36.2 million) but these exports plummeted to US$3.5 million the following year, according to the Comtrade database of the United Nations Statistics Division.[150]
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+ The Caribbean Industrial Research Institute in Trinidad and Tobago facilitates climate change research and provides industrial support for R&D related to food security. It also carries out equipment testing and calibration for major industries.[150]
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+ The Caribbean Epidemiology Centre in Port of Spain, University of Trinidad and Tobago, Tobago Institute of Health, and University of the West Indies (St Augustine campus) also conduct R&D.[150]
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+ Probability of a hurricane striking Caribbean countries in a given year, 2012 (%).[154]
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+ Electricity costs for the CARICOM countries, 2011[155]
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+ GERD by sector of performance in Trinidad and Tobago, 2000–2012[156]
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+ Scientific publication trends in the CARICOM countries, 2005–2014[157]
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+ Scientific publications in the CARICOM countries, 2014[157]
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+ USPTO patents granted to Caribbean countries, 2008–2013.[158]
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+ Trinidad and Tobago has a diverse culture mixing Indian, African, Creole, Chinese, Amerindian, Arab, Latino, and European influences, reflecting the various communities who have migrated to the islands over the centuries. The island is particularly renowned for its annual Carnival celebrations.[16] Festivals rooted in various religions and cultures practiced on the islands are also popular, such as Christmas, Divali, Phagwah (Holi), Easter, New Year’s Day, Hosay, Eid al-Fitr, the Santa Rosa Indigenous Festival, and Chinese New Year.
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+ Trinidad and Tobago claims two Nobel Prize-winning authors, V. S. Naipaul and St Lucian-born Derek Walcott (who also founded the Trinidad Theatre Workshop). Other notable writers include Neil Bissoondath, Vahni Capildeo, Earl Lovelace, Seepersad Naipaul, Shiva Naipaul, Lakshmi Persaud, Kenneth Ramchand, Arnold Rampersad, and Samuel Selvon.
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+ Trinidadian designer Peter Minshall is renowned not only for his Carnival costumes but also for his role in opening ceremonies of the Barcelona Olympics, the 1994 FIFA World Cup, the 1996 Summer Olympics, and the 2002 Winter Olympics, for which he won an Emmy Award.[161]
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+ Trinidad and Tobago is the birthplace of calypso music and the steelpan.[162][163][164] Trinidad is also the birthplace of soca music, chutney music, chutney-soca, parang, rapso, pichakaree and chutney parang.
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+ The limbo dance originated in Trinidad as an event that took place at wakes in Trinidad. The limbo has African roots. It was popularized in the 1950s by dance pioneer Julia Edwards[165] (known as the First Lady of Limbo) and her company which appeared in several films.[166] Bélé, Bongo, and whining are also dance forms with African roots.[167]
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+ Jazz, ballroom, ballet, modern, and salsa dancing are also popular.[167]
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+ Indian dance forms are also popular in Trinidad and Tobago.[168] Kathak, Odissi, and Bharatanatyam are the most popular Indian classical dance forms in Trinidad and Tobago.[169] Indian folk dances and Bollywood dances are also popular.[169]
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+ Geoffrey Holder (brother of Boscoe Holder) and Heather Headley are two Trinidad-born artists who have won Tony Awards for theatre. Holder also has a distinguished film career, and Headley has won a Grammy Award as well.
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+ Indian theatre is also popular throughout Trinidad and Tobago. Dramas such as Nautankis, Raja Harishchandra, Sharwan Kumar, and Alha-Khand were brought by Indians to Trinidad and Tobago, however they had largely began to die out, till preservation began by Indian cultural groups.[170] The drama about the life of the Hindu god Rama, Ramleela, is popular during the time between Sharad Navaratri and Dushera and the drama about the life of the Hindu god Krishna, Ras leela (Krishna leela), is popular around the time of Krishna Janmashtami.[171][172][173]
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+ Trinidad and Tobago is also smallest country to have two Miss Universe titleholders and the first black woman ever to win: Janelle Commissiong in 1977, followed by Wendy Fitzwilliam in 1998; the country has also had one Miss World titleholder, Giselle LaRonde.[citation needed]
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+ Hasely Crawford won the first Olympic gold medal for Trinidad and Tobago in the men's 100-metre dash in the 1976 Summer Olympics. Nine different athletes from Trinidad and Tobago have won twelve medals at the Olympics, beginning with a silver medal in weightlifting, won by Rodney Wilkes in 1948,[174] and most recently, a gold medal by Keshorn Walcott in the men's javelin throw in 2012. Ato Boldon has won the most Olympic and World Championship medals for Trinidad and Tobago in athletics, with eight in total – four from the Olympics and four from the World Championships. Boldon was the sole world champion Trinidad and Tobago has produced until Jehue Gordon in Moscow 2013. Ato won the 1997 200 m sprint World Championship in Athens. Swimmer George Bovell III won a bronze medal in the men's 200 m IM in 2004. At the 2017 World Championship in London, the Men 4x400 relay team captured the title, thus the country now celebrates three world championships titles. The team consisted of Jarrin Solomon, Jareem Richards, Machel Cedenio and Lalonde Gordon with Renny Quow who ran in the heats.
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+ Also in 2012 Lalonde Gordon competed in the XXX Summer Olympics where he won a bronze medal in the 400 metres (1,300 feet), being surpassed by Luguelin Santos of the Dominican Republic and Kirani James of Grenada. Keshorn Walcott (as stated above) came first in javelin and earned a gold medal, making him the second Trinidadian in the country's history to receive one. This also makes him the first Western[clarification needed] athlete in 40 years to receive a gold medal in the javelin sport, and the first athlete from Trinidad and Tobago to win a gold medal in a field event in the Olympics. Sprinter Richard Thompson is also from Trinidad and Tobago. He came second place to Usain Bolt in the Beijing Olympics in the 100 metres (330 feet) with a time of 9.89s.
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+ In 2018 The Court of Arbitration for Sport made its final decision on the failed doping sample from the Jamaican team in the 4 x 100 relay in the 2008 Olympic Games. The team from Trinidad and Tobago will be awarded the gold medal, because of the second rank during the relay run.[175]
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+ Cricket is a popular sport of Trinidad and Tobago, often deemed the national sport, and there is intense inter-island rivalry with its Caribbean neighbours. Trinidad and Tobago is represented at Test cricket, One Day International as well as Twenty20 cricket level as a member of the West Indies team. The national team plays at the first-class level in regional competitions such as the Regional Four Day Competition and Regional Super50. Meanwhile, the Trinbago Knight Riders play in the Caribbean Premier League.
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+ The Queen's Park Oval located in Port of Spain is the largest cricket ground in the West Indies, having hosted 60 Test matches as of January 2018. Trinidad and Tobago along with other islands from the Caribbean co-hosted the 2007 Cricket World Cup.
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+ Brian Lara, world record holder for the most runs scored both in a Test and in a First Class innings and other records, was born in a small town of Santa Cruz and is often referred to as the Prince of Port of Spain or simply the Prince. This legendary West Indian batsman is widely regarded (along with Sir Donald Bradman, Sunil Gavaskar and Sachin Tendulkar[citation needed]) as one of the best batsmen ever to have played the game,[citation needed] and is one of the most famous sporting icons in the country.
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+ Association football is also a popular sport in Trinidad and Tobago. The men's national football team qualified for the 2006 FIFA World Cup for the first time by beating Bahrain in Manama on 16 November 2005, making them the second smallest country ever (in terms of population) to qualify, after Iceland. The team, coached by Dutchman Leo Beenhakker, and led by Tobagonian-born captain Dwight Yorke, drew their first group game – against Sweden in Dortmund, 0–0, but lost the second game to England on late goals, 0–2. They were eliminated after losing 2–0 to Paraguay in the last game of the Group stage. Prior to the 2006 World Cup qualification, Trinidad and Tobago came close in a controversial qualification campaign for the 1974 FIFA World Cup. Following the match, the referee of their critical game against Haiti was awarded a lifetime ban for his actions.[176] Trinidad and Tobago again fell just short of qualifying for the World Cup in 1990, needing only a draw at home against the United States but losing 1–0.[177] They play their home matches at the Hasely Crawford Stadium. Trinidad and Tobago hosted the 2001 FIFA U-17 World Championship, and hosted the 2010 FIFA U-17 Women's World Cup.
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+ The TT Pro League is the country's primary football competition and is the top level of the Trinidad and Tobago football league system. The Pro League serves as a league for professional football clubs in Trinidad and Tobago. The league began in 1999 as part of a need for a professional league to strengthen the country's national team and improve the development of domestic players. The first season took place in the same year beginning with eight teams.
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+ Basketball is commonly played in Trinidad and Tobago in colleges, universities and throughout various urban basketball courts. Its national team is one of the most successful teams in the Caribbean. At the Caribbean Basketball Championship it won four straight gold medals from 1986 to 1990.
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+ Netball has long been a popular sport in Trinidad and Tobago, although it has declined in popularity in recent years. At the Netball World Championships they co-won the event in 1979, were runners up in 1987, and second runners up in 1983.
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+ Rugby is played in Trinidad and Tobago and continues to be a popular sport, and horse racing is regularly followed in the country.
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+ There is also the Trinidad and Tobago national baseball team which is controlled by the Baseball/Softball Association of Trinidad and Tobago, and represents the nation in international competitions. The team is a provisional member of the Pan American Baseball Confederation.
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+ There are a number of 9 and 18-hole golf courses on Trinidad and Tobago. The most established is the St Andrews Golf Club, Maraval in Trinidad (commonly referred to as Moka), and there is a newer course at Trincity, near Piarco Airport called Millennium Lakes. There are 18-hole courses at Chaguramas and Point-a-Pierre and 9-hole courses at Couva and St Madeline. Tobago has two 18-hole courses. The older of the two is at Mount Irvine, with the Magdalena Hotel & Golf Club (formerly Tobago Plantations) being built more recently.
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+ Although a minor sport, bodybuilding is of growing interest in Trinidad and Tobago. Heavyweight female bodybuilder Kashma Maharaj is of Trinidadian descent.
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+ Dragonboat is also another water-sport that has been rapidly growing over the years. Introduced in 2006. the fraternity made consistent strides in having more members apart of the TTDBF (Trindad and Tobago Dragonboat Federation) as well as performing on an international level such as the 10th IDBF World Nations Dragon Boat Championships in Tampa, Florida in the US in 2011.
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+ Claude Noel is a former world champion in professional boxing. He was born in Tobago.
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+ The flag was chosen by the Independence committee in 1962. Red, black and white symbolise the warmth of the people, the richness of the earth and water respectively.[178][179]
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+ The coat of arms was designed by the Independence committee, and features the scarlet ibis (native to Trinidad), the cocrico (native to Tobago) and hummingbird. The shield bears three ships, representing both the Trinity, and the three ships that Columbus sailed.[178]
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+ There are five categories and thirteen classes of national awards:[180]
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+ The national anthem of the twin-island state is "Forged from the Love of Liberty".[181][182]
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+ Other national songs include "God Bless Our Nation"[183] and "Our Nation's Dawning".[184]
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+
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+ The national flower of Trinidad and Tobago is the chaconia flower. It was chosen as the national flower because it is an indigenous flower that has witnessed the history of Trinidad and Tobago. It was also chosen as the national flower because of its red colour that resembles the red of the national flag and coat of arms and because it blooms around the Independence Day of Trinidad and Tobago.[185]
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+
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+ The national birds of Trinidad and Tobago are the scarlet ibis and the cocrico. The scarlet ibis is kept safe by the government by living in the Caroni Bird Sanctuary which was set up by the government for the protection of these birds. The Cocrico is more indigenous to the island of Tobago and are more likely to be seen in the forest.[186] The hummingbird is considered another symbol of Trinidad and Tobago due to its significance to the indigenous peoples, however, it is not a national bird.[187]
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+ The scarlet ibis birds flying over the Caroni Swamp
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+ The cocrico bird in Tobago
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+
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+
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+
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+ This article incorporates text from a free content work. UNESCO Science Report: towards 2030, 156–173, Harold Ramkissoon & Ishenkumba A. Kahwa, UNESCO Publishing. To learn how to add open license text to Wikipedia articles, please see this how-to page. For information on reusing text from Wikipedia, please see the terms of use.
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+ ^ These three Dutch Caribbean territories form the SSS islands.
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+ * These three Dutch Caribbean territories form the BES islands.
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+ † Physiographically, these are continental islands not a part of the volcanic Windward Islands arc. However, based on proximity, these islands are sometimes grouped with the Windward Islands culturally and politically.
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+
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+ ~ Disputed territories administered by Colombia.
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+ # Physiographically, Bermuda is an isolated oceanic island in the North Atlantic Ocean, not a part of the Antilles, West Indies, Caribbean, North American continent or South American continent. Usually grouped with Northern American countries based on proximity; occasionally grouped with the Caribbean region culturally.
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+ Coordinates: 10°36′N 61°6′W / 10.600°N 61.100°W / 10.600; -61.100
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1
+
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+
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+ Trinidad and Tobago (/ˈtrɪnɪdæd ... təˈbeɪɡoʊ/ (listen), /- toʊ-/), officially the Republic of Trinidad and Tobago, is the southernmost island country in the Caribbean.[14][15] Consisting of the main islands Trinidad and Tobago, and numerous much smaller islands, it is situated 130 kilometres (81 miles) south of Grenada and 11 kilometres (6.8 miles) off the coast of northeastern Venezuela.[16] It shares maritime boundaries with Barbados to the northeast, Grenada to the northwest, Guyana to the southeast, and Venezuela to the south and west.[17][18]
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+
5
+ The island of Trinidad was inhabited for centuries by native Amerindian peoples before becoming a colony in the Spanish Empire, following the arrival of Christopher Columbus in 1498. Spanish governor Don José María Chacón surrendered the island to a British fleet under the command of Sir Ralph Abercromby in 1797.[19] During the same period, the island of Tobago changed hands among Spanish, British, French, Dutch and Courlander colonisers more times than any other island in the Caribbean.[citation needed] Trinidad and Tobago were ceded to Britain in 1802 under the Treaty of Amiens as separate states and unified in 1889.[20] Trinidad and Tobago obtained independence in 1962, becoming a republic in 1976.[15][16]
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+ Trinidad and Tobago has the third highest GDP per capita based on purchasing power parity (PPP) in the Americas after the United States and Canada.[21] It is recognised by the World Bank as a high-income economy.[22] Unlike most Caribbean nations and territories, which rely heavily on tourism, the Trinidadian economy is primarily industrial with an emphasis on petroleum and petrochemicals;[23] much of the nation's wealth is derived from its large reserves of oil and natural gas.[24][25]
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+ Trinidad and Tobago is well known for its African and Indian cultures, reflected in its large and famous Carnival, Diwali, and Hosay celebrations, as well being the birthplace of steelpan, the limbo, and music styles such as calypso, soca, rapso, parang, chutney, and chutney soca.[26][27][28][29][30][31][32][33]
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+
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+ Historian E. L. Joseph claimed that Trinidad's Amerindian name was Cairi or "Land of the Humming Bird", derived from the Arawak name for hummingbird, ierèttê or yerettê.[citation needed] However, other authors dispute this etymology with some claiming that cairi does not mean hummingbird (tukusi or tucuchi being suggested as the correct word) and some claiming that kairi, or iere, simply means island.[clarification needed][34][citation needed] Christopher Columbus renamed it "La Isla de la Trinidad" ("The Island of the Trinity"), fulfilling a vow made before setting out on his third voyage of exploration.[35] Tobago's cigar-like shape, or the use of tobacco by the native people, may have given it its Spanish name (cabaco, tavaco, tobacco) and possibly some of its other Amerindian names, such as Aloubaéra (black conch) and Urupaina (big snail),[34] although the English pronunciation is /təˈbeɪɡoʊ/.[citation needed]
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+
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+ Both Trinidad and Tobago were originally settled by Amerindians who came through South America.[16] Trinidad was first settled by pre-agricultural Archaic people at least 7,000 years ago, making it the earliest settled part of the Caribbean.[36] Banwari Trace in south-west Trinidad is the oldest attested archaeological site in the Caribbean, dating to about 5000 BC. Several waves of migration occurred over the following centuries, which can be identified by differences in their archaeological remains.[37] At the time of European contact, Trinidad was occupied by various Arawakan-speaking groups including the Nepoya and Suppoya, and Cariban-speaking groups such as the Yao, while Tobago was occupied by the Island Caribs and Galibi. Trinidad was known to the native peoples as 'Ieri' ('Land of the Humming Bird').[36]
14
+
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+ Christopher Columbus was the first European to see Trinidad, on his third voyage to the Americas in 1498.[36][38] He also reported seeing Tobago on the distant horizon, naming it Bellaforma, but did not land on the island.[16][39]
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+
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+ In the 1530s Antonio de Sedeño, a Spanish soldier intent on conquering the island of Trinidad, landed on its southwest coast with a small army of men, intending to subdue the Amerindian peoples of the island. Sedeño and his men fought the native peoples on many occasions, and subsequently built a fort. The next few decades were generally spent in warfare with the native peoples, until in 1592, the 'Cacique' (native chief) Wannawanare (also known as Guanaguanare) granted the area around modern Saint Joseph to Domingo de Vera e Ibargüen, and withdrew to another part of the island.[34] The settlement of San José de Oruña was later established by Antonio de Berrío on this land in 1592.[16][36] Shortly thereafter the English sailor Sir Walter Raleigh arrived in Trinidad on 22 March 1595 in search of the long-rumoured "El Dorado" ('City of Gold') supposedly located in South America.[36] He attacked San José, captured and interrogated Antonio de Berrío, and obtained much information from him and from the Cacique Topiawari; Raleigh then went on his way, and Spanish authority was restored.[40]
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+
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+ Meanwhile, there were numerous attempts by European powers to settle Tobago during the 1620-40s, with the Dutch, English and Courlanders (people from the Duchy of Courland and Semigallia, now part of Latvia) all attempting to colonise the island with little success.[41][42] From 1654 the Dutch and Courlanders managed to gain a more secure foothold, later joined by several hundred French settlers.[41] A plantation economy developed based on the production of sugar, indigo and rum, worked by large numbers of African slaves who soon came to vastly outnumber the European colonists.[42][41] Large numbers of forts were constructed as Tobago became a source of contention between France, Holland and Britain, with the island changing hands some 31 times prior to 1814, a situation exacerbated by widespread piracy.[42] The British managed to hold Tobago from 1762–1781, whereupon it was captured by the French, who ruled until 1793 when Britain re-captured the island.[42]
20
+
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+ The 17th century on Trinidad passed largely without major incident, but sustained attempts by the Spaniards to control and rule over the Amerindians were often fiercely resisted.[36] In 1687 the Catholic Catalan Capuchin friars were given responsibility for the conversions of the indigenous people of Trinidad and the Guianas.[36] They founded several missions in Trinidad, supported and richly funded by the state, which also granted encomienda right to them over the native peoples, in which the native peoples were forced to provide labour for the Spanish.[36] One such mission was Santa Rosa de Arima, established in 1789, when Amerindians from the former encomiendas of Tacarigua and Arauca (Arouca) were relocated further west.[citation needed] Escalating tensions between the Spaniards and Amerindians culminated in violence 1689, when Amerindians in the San Rafael encomienda rebelled and killed several priests, attacked a church, and killed the Spanish governor José de León y Echales. Among those killed in the governor's party was Juan Mazien de Sotomayor, missionary priest to the Nepuyo villages of Cuara, Tacarigua and Arauca.[citation needed] The Spanish retaliated severely, slaughtering hundreds of native peoples in an event that became known as the Arena massacre.[36] As a result of this, continuing Spanish slave-raiding, and the devastating impact of introduced disease to which they had no immunity, the native population was virtually wiped out by the end of the following century.[43][36]
22
+
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+ During this period Trinidad was an island province belonging to the Viceroyalty of New Spain, together with Central America, present-day Mexico and the southwestern United States.[44] In 1757 the capital was moved from San José de Oruña to Puerto de España (modern Port of Spain) following several pirate attacks.[45] However the Spanish never made any concerted effort to colonise the islands; Trinidad in this period was still mostly forest, populated by a few Spaniards with a handful of slaves and a few thousand Amerindians.[44] Indeed, the population in 1777 was only 1,400, and Spanish colonisation in Trinidad remained tenuous.[citation needed]
24
+
25
+ Since Trinidad was considered underpopulated, Roume de St. Laurent, a Frenchman living in Grenada, was able to obtain a Cédula de Población from the Spanish king Charles III on 4 November 1783.[46] A Cédula de Población had previously been granted in 1776 by the king, but had not shown results, and therefore the new Cédula was more generous.[16] It granted free land and tax exemption for 10 years to Roman Catholic foreign settlers who were willing to swear allegiance to the King of Spain.[16] The Spanish also gave many incentives to lure settlers to the island, including exemption from taxes for ten years and land grants in accordance with the terms set out in the Cédula.[47] The land grant was 30 fanegas (13 hectares/32 acres) for each free man, woman and child and half of that for each slave that they brought with them. The Spanish sent a new governor, José María Chacón, to implement the terms of the new cédula.[46]
26
+
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+ It was fortuitous that the Cédula was issued only a few years before the French Revolution. During that period of upheaval, French planters with their slaves, free coloureds and mulattos from the neighbouring islands of Martinique, Saint Lucia, Grenada, Guadeloupe and Dominica migrated to Trinidad, where they established an agriculture-based economy (sugar and cocoa).[44] These new immigrants established local communities in Blanchisseuse, Champs Fleurs, Paramin,[48] Cascade, Carenage and Laventille.
28
+
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+ As a result, Trinidad's population jumped to over 15,000 by the end of 1789, and by 1797 the population of Port of Spain had increased from under 3,000 to 10,422 in just five years, with a varied population of mixed race individuals, Spaniards, Africans, French republican soldiers, retired pirates and French nobility.[44] The total population of Trinidad was 17,718, of which 2,151 were of European ancestry, 4,476 were "free blacks and people of colour", 10,009 were enslaved people and 1,082 Amerindians.[citation needed] The sparse settlement and slow rate of population-increase during Spanish rule (and even later during British rule) made Trinidad one of the less populated colonies of the West Indies, with the least developed plantation infrastructure.[49]
30
+
31
+ The British had begun to take a keen interest in Trinidad, and in 1797 a British force led by General Sir Ralph Abercromby launched an invasion of Trinidad.[16][50] His squadron sailed through the Bocas and anchored off the coast of Chaguaramas. Seriously outnumbered, Governor Chacón decided to capitulate to British without fighting.[50] Trinidad thus became a British crown colony, with a largely French-speaking population and Spanish laws.[44] British rule was later formalised under the Treaty of Amiens (1802).[16][50] The colony's first British governor was Thomas Picton, however his heavy-handed approach to enforcing British authority, including the use of torture and arbitrary arrest, led to his being recalled.[50]
32
+
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+ British rule led to an influx of settlers from the United Kingdom and the British colonies of the Eastern Caribbean. English, Scots, Irish, German and Italian families arrived, as well as some free blacks known as 'Merikins' who had fought for Britain in the War of 1812 and were granted land in southern Trinidad.[51][52][53] Under British rule, new states were created and the importation of slaves increased, however by this time support for abolitionism had vastly increased and in England the slave trade was under attack.[49][54] Slavery was abolished in 1833, after which former slaves served an "apprenticeship" period. In 1837 Daaga, a West African slave trader who had been captured by Portuguese slavers and later rescued by the British navy, was conscripted into the local regiment. Daaga and a group of his compatriots mutinied at the barracks in St Joseph and set out eastward in an attempt to return to their homeland. The mutineers were ambushed by a militia unit just outside the town of Arima. The revolt was crushed at the cost of some 40 dead, and Daaga and his party were later executed at St Joseph.[55] The apprenticeship system ended on 1 August 1838 with full emancipation.[16][53] An overview of the populations statistics in 1838, however, clearly reveals the contrast between Trinidad and its neighbouring islands: upon emancipation of the slaves in 1838, Trinidad had only 17,439 slaves, with 80% of slave owners having enslaved fewer than 10 people each.[56] In contrast, at twice the size of Trinidad, Jamaica had roughly 360,000 slaves.[57]
34
+
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+ After the African slaves were emancipated many refused to continue working on the plantations, often moving out to urban areas such as Laventille and Belmont to the east of Port of Spain.[53] As a result, a severe agricultural labour shortage emerged; the British filled this gap by instituting a system of indentureship. Various nationalities were contracted under this system, including Indians, Chinese, and Portuguese.[58] Of these, the East Indians were imported in the largest numbers, starting from 1 May 1845, when 225 Indians were brought in the first shipment to Trinidad on the Fatel Razack, a Muslim-owned vessel.[53][59] Indentureship of the Indians lasted from 1845 to 1917, during which time more than 147,000 Indians came to Trinidad to work on sugarcane plantations.[16][60]
36
+
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+ Indentureship contracts were sometimes exploitative, to such an extent that historians such as Hugh Tinker were to call it "a new system of slavery". Despite these descriptions, it was not truly a new form of slavery, as workers were paid, contracts were finite, and the idea of an individual being another's property had been eliminated when slavery was abolished.[61] In addition, employers of indentured labour had no legal right to flog or whip their workers; the main legal sanction for the enforcement of the indenture laws was prosecution in the courts, followed by fines or (more likely) jail sentences.[62] People were contracted for a period of five years, with a daily wage as low as 25 cents in the early 20th century, and they were guaranteed return passage to India at the end of their contract period. However, coercive means were often used to retain labourers, and the indentureship contracts were soon extended to 10 years from 1854 after the planters complained that they were losing their labour too early.[49][53] In lieu of the return passage, the British authorities soon began offering portions of land to encourage settlement, and by 1902, more than half of the sugar cane in Trinidad was being produced by independent cane farmers; the majority of which were Indians.[63] Despite the trying conditions experienced under the indenture system, about 90% of the Indian immigrants chose, at the end of their contracted periods of indenture, to make Trinidad their permanent home.[63] East Indians entering the colony were also subject to certain crown laws which segregated them from the rest of Trinidad's population, such as the requirement that they carry a pass with them if they left the plantations, and that if freed, they carry their "Free Papers" or certificate indicating completion of the indenture period.[64]
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+
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+ Few Indians settled on Tobago however, and the descendants of African slaves continued to form the majority of the island's population. An ongoing economic slump in the middle-to-late 19th century caused
40
+ widespread poverty.[65] Discontent erupted into rioting on the Roxborough plantation in 1876, in an event known as the Belmanna Uprising after a policeman who was killed.[65] The British eventually managed to restore control, however as a result of the disturbances Tobago's Legislative Assembly voted to dissolve itself and the island became a Crown colony in 1877.[65] With the sugar industry in a state of near-collapse and the island no longer profitable, the British attached Tobago to their Trinidad colony in 1899.[16][66][67]
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+
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+ In 1903, a protest against the introduction of new water rates in Port of Spain erupted into rioting; 18 people were shot dead, and the Red House (the government headquarters) was damaged by fire.[66] A local elected assembly with some limited powers was introduced in 1913.[66] Economically Trinidad and Tobago remained a predominantly agricultural colony; alongside sugarcane, the cacao (cocoa) crop also contributed greatly to economic earnings in the late 19th and early 20th centuries.
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+
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+ In November 1919, the dockworkers went on strike over bad management practises, low wages compared to a higher cost of living.[68] Strikebreakers were brought in to keep a minimum of goods moving through the ports. On December 1, 1919, the striking dockworkers rushed the harbour and chased off the strikebreakers.[68] They then proceeded to march on the government buildings in Port of Spain. Other unions and workers, many with the same grievances, joined the dock worker's strike making it a General Strike.[68] Violence broke out and was only put down with help from the sailors of British Naval ship HMS Calcutta. The unity brought upon by the strike was the first time of cooperation between the various ethnic groups of the time.[69] Historian Brinsley Samaroo says that the 1919 strikes "seem to indicate that there was a growing class consciousness after the war and this transcended racial feelings at times."[69]
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+
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+ However, in the 1920s, the collapse of the sugarcane industry, concomitant with the failure of the cocoa industry, resulted in widespread depression among the rural and agricultural workers in Trinidad, and encouraged the rise of a labour movement. Conditions on the islands worsened in the 1930s with the onset of the Great Depression, with an outbreak of labour riots occurring in 1937 which resulted in several deaths.[70] The labour movement aimed to unite the urban working class and agricultural labour class; the key figures being Arthur Cipriani, who led the Trinidad Workingmen's Association (TWA), and Tubal Uriah "Buzz" Butler of the British Empire Citizens' and Workers' Home Rule Party.[70] As the movement developed calls for greater autonomy from British colonial rule became widespread; this effort was severely undermined by the British Home Office and by the British-educated Trinidadian elite, many of whom were descended from the plantocracy class.
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+
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+ Petroleum had been discovered in 1857, but became economically significant only in the 1930s and afterwards as a result of the collapse of sugarcane and cocoa, and increasing industrialisation.[71]
49
+ [72][73] By the 1950s petroleum had become a staple in Trinidad's export market, and was responsible for a growing middle class among all sections of the Trinidad population. The collapse of Trinidad's major agricultural commodities, followed by the Depression, and the rise of the oil economy, led to major changes in the country's social structure.
50
+
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+ The presence of American military bases in Chaguaramas and Cumuto in Trinidad during World War II had a profound effect on society. The Americans vastly improved the infrastructure on Trinidad and provided many locals with well-paying jobs; however the social effects of having so many young soldiers stationed on the island, as well as their often unconcealed racial prejudice, caused resentment.[66] The Americans left in 1961.[74]
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+
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+ In the post-war period the British began a process of decolonisation across the British Empire. In 1945 universal suffrage was introduced to Trinidad and Tobago.[16][66] Political parties emerged on the island, however these were largely divided along racial lines: Afro-Trinidadians and Tobagonians primarily supported the People's National Movement (PNM), formed in 1956 by Eric Williams, with Indo-Trinidadians and Tobagonians mostly supporting the People's Democratic Party (PDP), formed in 1953 by Bhadase Sagan Maraj,[75] which later merged into the Democratic Labour Party (DLP) in 1957.[76] Britain's Caribbean colonies formed the West Indies Federation in 1958 as a vehicle for independence, however the Federation dissolved after Jamaica withdrew following a membership referendum in 1961. The government of Trinidad and Tobago subsequently chose to seek independence from the United Kingdom on its own.[77]
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+
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+ Trinidad and Tobago gained its independence from the United Kingdom on 31 August 1962.[16][73] Elizabeth II remained head of state as Queen of Trinidad and Tobago, represented locally by Governor-General Solomon Hochoy. Eric Williams of the PNM, a noted historian and intellectual widely regarded as The Father of The Nation, became the first Prime Minister, serving in that capacity uninterrupted until 1981.[16] The dominant figure in the opposition in the early independence years was Rudranath Capildeo of the DLP. The 1960s saw the rise of a Black Power movement, inspired in part by the civil rights movement in the United States. Protests and strikes became common, with events coming to head in April 1970 when police shot dead a protester named Basil Davis.[76] Fearing a breakdown of law and order, Prime Minister Williams declared a state of emergency and arrested many of the Black Power leaders. Some army leaders who were sympathetic to the Black Power movement, notably Raffique Shah and Rex Lassalle, attempted to mutiny; however, this was quashed by the Trinidad and Tobago Coast Guard.[76] Williams and the PNM retained power, largely due to divisions in the opposition.[76]
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+
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+ In 1963 Tobago was struck by Hurricane Flora, which killed 30 people and resulted in enormous destruction across the island.[78] Partly as a result of this, tourism came to replace agriculture as the island's main income earner in the subsequent decades.[78]
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+
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+ Between the years 1972 and 1983, the country profited greatly from the rising price of oil and the discovery of vast new oil deposits in its territorial waters, resulting in an economic boom that increased living standards greatly.[16][76] In 1976 the country became a republic within the Commonwealth, though it retained the Judicial Committee of the Privy Council as its final appellate court.[16] The position of governor-general was replaced with that of President; Ellis Clarke was the first to hold this largely ceremonial role.[79] Tobago was granted limited self-rule with the creation of the Tobago House of Assembly in 1980.[65]
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+
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+ Williams died in 1981, being replaced by George Chambers who led the country until 1986. By this time a fall in the price of oil had resulted in a recession, causing rising inflation and unemployment.[80] The main opposition parties united under the banner of National Alliance for Reconstruction (NAR) and won the 1986 Trinidad and Tobago general election, with NAR leader A. N. R. Robinson becoming the new Prime Minister.[81][76] Robinson was unable to hold together the fragile NAR coalition, and social unrest was caused by his economic reforms, such as devaluing the currency and implementing an International Monetary Fund Structural Adjustment Program.[16] In 1990 114 members of the Jamaat al Muslimeen, led by Yasin Abu Bakr (formerly known as Lennox Phillip) stormed the Red House (the seat of Parliament), and Trinidad and Tobago Television, the only television station in the country at the time, holding Robinson and country's government hostage for six days before surrendering.[82] The coup leaders were promised amnesty, but upon their surrender they were then arrested, but later released after protracted legal wrangling.[58]
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+
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+ The PNM under Patrick Manning returned to power following the 1991 Trinidad and Tobago general election.[16] Hoping to capitalise on an improvement in the economy, Manning called an early election in 1995, however, this resulted in a hung parliament. Two NAR representatives backed the opposition United National Congress (UNC), which had split off from the NAR in 1989, and they thus took power under Basdeo Panday, who became the country's first Indo-Trinidadian Prime Minister.[16][80][83] After a period of political confusion caused by a series of inconclusive election results, Patrick Manning returned to power in 2001, retaining that position until 2010.[16]
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+
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+ Since 2003 the country entered a second oil boom, and petroleum, petrochemicals and natural gas continue to be the backbone of the economy. Tourism and the public service are the mainstay of the economy of Tobago, though authorities have attempted to diversify the island's economy.[84] A corruption scandal resulted in Manning's defeat by the newly formed People's Partnership coalition in 2010, with Kamla Persad-Bissessar becoming the country's first female Prime Minister.[85][86][87] However, corruption allegations bedevilled the new administration, and the PP were defeated in 2015 by the PNM under Keith Rowley.[88][89]
66
+
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+ Trinidad and Tobago is situated between 10° 2' and 11° 12' N latitude and 60° 30' and 61° 56' W longitude, with the Caribbean Sea to the north, the Atlantic Ocean to the east and south, and the Gulf of Paria to the west. It is located in the far south-east of the Caribbean region, with the island of Trinidad being just 11 kilometres (6.8 mi) off the coast of Venezuela in mainland South America across the Columbus Channel.[16] Covering an area of 5,128 km2 (1,980 sq mi),[90] the country consists of two main islands, Trinidad and Tobago, separated by a 20m (30 km) strait, plus a number of much smaller islands, including Chacachacare, Monos, Huevos, Gaspar Grande (or Gasparee), Little Tobago, and Saint Giles Island.[16]
68
+
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+ Trinidad is 4,768 km2 (1,841 sq mi) in area (comprising 93.0% of the country's total area) with an average length of 80 kilometres (50 mi) and an average width of 59 kilometres (37 mi). Tobago has an area of about 300 km2 (120 sq mi), or 5.8% of the country's area, is 41 km (25 mi) long and 12 km (7.5 mi) at its greatest width. Trinidad and Tobago lie on the continental shelf of South America, and are thus geologically considered to lie entirely in South America.[16]
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+
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+ The terrain of the islands is a mixture of mountains and plains.[15] On Trinidad the Northern Range runs parallel with the north coast, and contains the country's highest peak (El Cerro del Aripo), which is 940 metres (3,080 ft) above sea level[15] and second highest (El Tucuche, 936 metres (3,071 ft)).[16] The rest of the island is generally flatter, excluding the Central Range and Montserrat Hills in the centre of the island and the Southern Range and Trinity Hills in the south. The east coast is noted for its beaches, most notably Manzanilla Beach. The island contains several large swamp areas, such as the Caroni Swamp and the Nariva Swamp.[16] Major bodies of water on Trinidad include the Hollis Reservoir, Navet Reservoir, Caroni Reservoir. Trinidad is made up of a variety of soil types, the majority being fine sands and heavy clays. The alluvial valleys of the Northern Range and the soils of the East–West Corridor are the most fertile.[91][citation needed] Trinidad is also notable for containing Pitch Lake, the largest natural reservoir of asphalt in the world.[15][16] Tobago contains a flat plain in its south-west, with the eastern half of the island being more mountainous, culminating in Pigeon Peak, the island's highest point at 550 metres (1,800 ft).[92] Tobago also contains several coral reefs off its coast.[16]
72
+
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+ The majority of the population reside on the island of Trinidad, and this is thus the location of largest towns and cities. There are four major municipalities in Trinidad: the capital Port of Spain, San Fernando, Arima and Chaguanas. The main town on Tobago is Scarborough.
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+
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+ The Northern Range consists mainly of Upper Jurassic and Cretaceous metamorphic rocks. The Northern Lowlands (the East–West Corridor and Caroni Plain) consist of younger shallow marine clastic sediments. South of this, the Central Range fold and thrust belt consists of Cretaceous and Eocene sedimentary rocks, with Miocene formations along the southern and eastern flanks. The Naparima Plain and the Nariva Swamp form the southern shoulder of this uplift.[citation needed]
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+
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+ The Southern Lowlands consist of Miocene and Pliocene sands, clays, and gravels. These overlie oil and natural gas deposits, especially north of the Los Bajos Fault. The Southern Range forms the third anticlinal uplift. The rocks consist of sandstones, shales, siltstones and clays formed in the Miocene and uplifted in the Pleistocene. Oil sands and mud volcanoes are especially common in this area.[citation needed]
78
+
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+ Trinidad and Tobago has a maritime tropical climate.[15][16] There are two seasons annually: the dry season for the first five months of the year, and the rainy season in the remaining seven of the year. Winds are predominantly from the northeast and are dominated by the northeast trade winds. Unlike many Caribbean islands Trinidad and Tobago lies outside the main hurricane alleys; nevertheless, the island of Tobago was struck by Hurricane Flora on September 30, 1963. In the Northern Range of Trinidad, the climate is often cooler than that of the sweltering heat of the plains below, due to constant cloud and mist cover, and heavy rains in the mountains.
80
+
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+ Record temperatures for Trinidad and Tobago are 39 °C (102 °F)[93] for the high in Port of Spain, and a low of 12 °C (54 °F).[94]
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+
83
+ Because Trinidad and Tobago lies on the continental shelf of South America, and in ancient times were physically connected to the South American mainland, its biological diversity is unlike that of most other Caribbean islands, and has much more in common with that of Venezuela.[95] The main ecosystems are: coastal and marine (coral reefs, mangrove swamps, open ocean and seagrass beds); forest; freshwater (rivers and streams); karst; man-made ecosystems (agricultural land, freshwater dams, secondary forest); and savanna. On 1 August 1996, Trinidad and Tobago ratified the 1992 Rio Convention on Biological Diversity, and it has produced a biodiversity action plan and four reports describing the country's contribution to biodiversity conservation. These reports formally acknowledged the importance of biodiversity to the well-being of the country's people through provision of ecosystem services.[96]
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+
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+ Information about vertebrates is good, with 472 bird species (2 endemics), about 100 mammals, about 90 reptiles (a few endemics), about 30 amphibians (including several endemics), 50 freshwater fish and at least 950 marine fish.[97] Notable mammal species include the ocelot, manatee, collared peccary (known as the quenk locally), agouti, lappe, red brocket deer, otter, weeper capuchin and red howler monkey; there are also some 70 species of bat, including the vampire bat and fringe-lipped bat.[16][98] Amongst the reptiles, the spectacled caiman is the largest, sometimes growing up to 3m.[95] There are also 47 species of snake, including four venomous species, lizards such as the gecko, iguana, matte lizard and also several species of turtle.[16][99] are present. Of the amphibians, the golden tree frog is endemic to Trinidad.[99] Marine life is abundant, with several species of sea urchin, coral, lobster, anemone, starfish, manta ray, dolphin, porpoise and whale shark present in the islands' waters.[100] The introduced lionfish is viewed as a pest, as it eats many native species of fish and has no natural predators; efforts are currently underway to cull the numbers of this species.[100]
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+
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+ Trinidad and Tobago is noted particularly for its large number of bird species, and is a popular destination for bird watchers. Notable species include the scarlet ibis, cocrico, egret, shiny cowbird, bananaquit, oilbird and various species of honeycreeper, trogon, toucan, parrot, tanager, woodpecker, antbird, kites, hawks, boobies, pelicans and vultures; there are also 17 species of hummingbird, including the tufted coquette which is the world's third smallest.[101]
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+
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+ Information about invertebrates is dispersed and very incomplete. About 650 butterflies,[97] at least 672 beetles (from Tobago alone)[102] and 40 corals[97] have been recorded.[97] Other notable invertebrates include the cockroach, leaf-cutter ant and numerous species of mosquitoes, termites, spiders and tarantulas.
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+
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+ Although the list is far from complete, 1,647 species of fungi, including lichens, have been recorded.[103][104][105] The true total number of fungi is likely to be far higher, given the generally accepted estimate that only about 7% of all fungi worldwide have so far been discovered.[106] A first effort to estimate the number of endemic fungi tentatively listed 407 species.[107]
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+
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+ Information about micro-organisms is dispersed and very incomplete. Nearly 200 species of marine algae have been recorded.[97] The true total number of micro-organism species must be much higher.
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+
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+ Thanks to a recently published checklist, plant diversity in Trinidad and Tobago is well documented with about 3,300 species (59 endemic) recorded.[97] Despite significant felling, forests still cover about 40% of the country, and there are about 350 different species of tree.[95] A notable tree is the manchineel which is extremely poisonous to humans, and even just touching its sap can cause severe blistering of the skin; the tree is often covered with warning signs.
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+
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+ Trinidad and Tobago is a republic with a two-party system and a bicameral parliamentary system based on the Westminster System.[15]
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+
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+ The head of state of Trinidad and Tobago is the President, currently Paula Mae Weekes.[15] This largely ceremonial role replaced that of the Governor-General (representing the Monarch of Trinidad and Tobago) upon Trinidad and Tobago's becoming a republic in 1976.[16] The head of government is the Prime Minister, currently Keith Rowley.[15] The President is elected by an Electoral college consisting of the full membership of both houses of Parliament.
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+ The Prime Minister is elected following a general election which takes place every five years. The President is required to appoint the leader of the party who in his or her opinion has the most support of the members of the House of Representatives to this post; this has generally been the leader of the party which won the most seats in the previous election (except in the case of the 2001 General Elections).[16]
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+
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+ Since 1980 Tobago has also had its own elections, separate from the general elections. In these elections, members are elected and serve in the unicameral Tobago House of Assembly.[108][15][16]
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+
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+ Parliament consists of the Senate (31 seats) and the House of Representatives (41 seats, plus the Speaker).[109][15] The members of the Senate are appointed by the president; 16 Government Senators are appointed on the advice of the Prime Minister, six Opposition Senators are appointed on the advice of the Leader of the Opposition, currently Kamla Persad-Bissessar, and nine Independent Senators are appointed by the President to represent other sectors of civil society. The 41 members of the House of Representatives are elected by the people for a maximum term of five years in a "first past the post" system.
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+
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+ Trinidad is split into 14 regional corporations and municipalities, consisting of nine regions and five municipalities, which have a limited level of autonomy.[15][16] The various councils are made up of a mixture of elected and appointed members. Elections are held every three years.[citation needed] The country was formerly divided into counties.
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+
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+ The two main parties are the People's National Movement (PNM) and the United National Congress (UNC); another recent party was the Congress of the People (COP). Support for these parties appears to fall along ethnic lines, with the PNM consistently obtaining a majority of Afro-Trinidadian vote, and the UNC gaining a majority of Indo-Trinidadian support.
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+
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+ The Trinidad and Tobago Defence Force (TTDF) is the military organisation responsible for the defence of the twin island Republic of Trinidad and Tobago.[15] It consists of the Regiment, the Coast Guard, the Air Guard and the Defence Force Reserves. Established in 1962 after Trinidad and Tobago's independence from the United Kingdom, the TTDF is one of the largest military forces in the Anglophone Caribbean.[citation needed]
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+
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+ Its mission statement is to "defend the sovereign good of The Republic of Trinidad and Tobago, contribute to the development of the national community and support the State in the fulfillment of its national and international objectives". The Defence Force has been engaged in domestic incidents, such as the 1990 Coup Attempt, and international missions, such as the United Nations Mission in Haiti between 1993 and 1996.
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+
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+ In 2019, Trinidad and Tobago signed the UN treaty on the Prohibition of Nuclear Weapons.[110]
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+
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+ Trinidad and Tobago maintains close relations with its Caribbean neighbours and major North American and European trading partners. As the most industrialised and second-largest country in the Anglophone Caribbean, Trinidad and Tobago has taken a leading role in the Caribbean Community (CARICOM), and strongly supports CARICOM economic integration efforts. It also is active in the Summit of the Americas process and supports the establishment of the Free Trade Area of the Americas, lobbying other nations for seating the Secretariat in Port of Spain.[citation needed]
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+
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+ As a member of CARICOM, Trinidad and Tobago strongly backed efforts by the United States to bring political stability to Haiti, contributing personnel to the Multinational Force in 1994. After its 1962 independence, Trinidad and Tobago joined the United Nations and Commonwealth of Nations. In 1967 it became the first Commonwealth country to join the Organization of American States (OAS).[111] In 1995 Trinidad played host to the inaugural meeting of the Association of Caribbean States and has become the seat of this 35-member grouping, which seeks to further economic progress and integration among its states. In international forums, Trinidad and Tobago has defined itself as having an independent voting record, but often supports US and EU positions.[citation needed]
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+
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+ Trinidad and Tobago has in recent decades suffered from a relatively high crime rate;[112][113] there are currently roughly 500 murders per year.[114][76] The country is a noted transshipment centre for the trafficking of illegal drugs from South America to the rest of the Caribbean and beyond to North America.[115] Some estimates put the size of the 'hidden economy' as high as 20–30% of measured GDP.[116]
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+
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+ Though there have been no terrorism-related incidents in the country since the 1990 Islamic coup attempt, Trinidad and Tobago remains a potential target; for example, in February 2018 a plan to attack the Carnival was foiled by police.[113] It is estimated that roughly 100 citizens of the country have traveled to the Middle East to fight for Islamic State.[112][113] In 2017 the government adopted a counter-terrorism and extremism strategy.[113]
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+
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+ The country's prison administration is the Trinidad and Tobago Prison Service (TTPS), it is under the control of the Commissioner of Prisons Gerard Wilson, located in Port-of-Spain.[117] The prison population rate is 292 people per 100,000. The total prison population, including pre-trial detainees and remand prisoners, is 3,999 prisoners. The population rate of pre-trial detainees and remand prisoners is 174 per 100,000 of the national population (59.7% of the prison population). In 2018, the female prison population rate is 8.5 per 100,000 of the national population (2.9% of the prison population). Prisoners that are minors makes up 1.9% of the prison population and foreigners prisoners make 0.8% of the prison population. The occupancy level of Trinidad and Tobago's prison system is at 81.8% capacity.[117] Trinidad and Tobago has nine prison establishments; Golden Grove Prison, Maximum Security Prison, Port of Spain Prison, Eastern Correctional Rehabilitation Centre, Remand Prison, Tobago Convict Prison, Carrera Convict Island Prison, Women's Prison and Youth Training and Rehabilitation Centre.[118] Trinidad and Tobago also use labor yards as prisons, or means of punishment.[119]
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+
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+ The population of the country currently stands at 1,363,985 (July 2019 est.).[citation needed]
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+
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+ The ethnic composition of Trinidad and Tobago reflects a history of conquest and immigration.[120] While the earliest inhabitants were of Amerindian heritage, the two dominant groups in the country are now those of South Asian and of African heritage. Indo-Trinidadian and Tobagonians make up the country's largest ethnic group (approximately 35.4%);[15] they are primarily the descendants of indentured workers from South Asia (mostly from India), brought to replace freed African slaves who refused to continue working on the sugar plantations. Through cultural preservation many residents of Indian descent continue to maintain traditions from their ancestral homeland. Indo-Trinidadians reside primarily on Trinidad; as of the 2011 census only 2.5% of Tobago's population was of Indian descent.[121]
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+
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+ Afro-Trinidadians and Tobagonians make up the country's second largest ethnic group, with approximately 34.2% of the population identifying as being of African descent.[15] The majority of people of an African background are the descendants of slaves forcibly transported to the islands from as early as the 16th century. This group constitute the majority on Tobago, at 85.2%.[121]
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+ The bulk of the rest of the population are those who identify as being of mixed heritage.[15] There are also small but significant minorities of people of Amerindian, European, Chinese, and Arab descent. Arima on Trinidad is a noted centre of Amerindian culture.[16]
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+
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+ English is the country's official language (the local variety of standard English is Trinidadian and Tobagonian English or more properly, Trinidad and Tobago Standard English, abbreviated as "TTSE"), but the main spoken language is either of two English-based creole languages (Trinidadian Creole or Tobagonian Creole), which reflects the Amerindian, European, African, and Asian heritage of the nation. Both creoles contain elements from a variety of African languages; Trinidadian English Creole, however, is also influenced by French and French Creole (Patois).[122]
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+
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+ A majority of the early Indian immigrants spoke the Bhojpuri and Awadhi dialect of Hindustani (Hindi-Urdu), which later formed into Trinidadian Hindustani (Hindi-Urdu), which became the lingua franca of Indo-Trinidadian and Tobagonians. From 1935 Indian films began showing to audiences in Trinidad; most of these were in the Standard Hindustani (Hindi-Urdu) dialect and this modified Trinidadian Hindustani slightly by adding Standard Hindu and Urdu phrases and vocabulary to Trinidadian Hindustani. Indian films also revitalised Hindustani among Indo-Trinidadian and Tobagonians.[123] Around the mid to late 1970s the lingua franca of Indo-Trinidaian and Tobagonians switched from Trinidadian Hindustani to a sort of Hindinised version of English. Today Hindustani survives on through Indo-Trinidadian and Tobagonian musical forms such as, Bhajan, Indian classical music, Indian folk music, Filmi, Pichakaree, Chutney, Chutney soca, and Chutney parang. Presently there are about 26,000 people, which is 5.53% of the Indo-Trinidadian and Tobagonian population, who speak Trinidadian Hindustani. Many Indo-Trinidadians and Tobagonians today speak a type of Hinglish that consist of Trinidadian and Tobagonian English that is heavily laced with Trinidadian Hindustani vocabulary and phrases and many Indo-Trinidadians and Tobagonians can recite phrases or prayers in Hindustani today. There are many places in Trinidad and Tobago that have names of Hindustani origin. Some phrases and vocabulary have even made their way into the mainstream English and English Creole dialects of the country.[124][125][126][127][128]
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+
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+ The Chinese language first came to Trinidad and Tobago in 1806, when the British had brought Chinese labourers in order to determine if they were fit to use as labourers after the abolition of slavery.[citation needed] About 2,645 Chinese immigrants arrived in Trinidad as indentured labour between 1853 and 1866.[citation needed] A majority of the people who immigrated in the 19th century were from southern China and spoke the Hakka and Yue dialects of Chinese. In the 20th century after the years of indentureship up to the present-day more Chinese people have immigrated to Trinidad and Tobago for business and they speak the dialects of the indenturees along with other Chinese dialects, such as Mandarin and Min.[125][129] J. Dyer Ball, writing in 1906, says: "In Trinidad there were, about twenty years ago, 4,000 or 5,000 Chinese, but they have decreased to probably about 2,000 or 3,000, [2,200 in 1900]. They used to work in sugar plantations, but are now principally shopkeepers, as well as general merchants, miners and railway builders,
139
+ etc."[130]
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+
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+ The indigenous languages were Yao on Trinidad and Karina on Tobago, both Cariban, and Shebaya on Trinidad, which was Arawakan.[125]
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+
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+ According to the 2011 census,[3] Roman Catholics were the largest single religious group in Trinidad and Tobago with 21.60% of the total population. The Pentecostal/Evangelical/Full Gospel denominations were the third largest group with 12.02% of the population. The remaining population is made of various Christian denominations (Spiritual Shouter Baptists (5.67%), Anglicans (5.67%), Seventh-day Adventists (4.09%), Presbyterians or Congregationalists (2.49%), Jehovah's Witnesses (1.47%), other Baptists (1.21%), Methodists (0.65%) and the Moravian Church (0.27%)). Respondents who did not state a religious affiliation represented 11.1% of the population, with 2.18% declaring themselves Irreligious.
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+
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+ Hindus were the second largest group with 18.15%.[3] Hinduism is practiced throughout the country and Diwali is a public holiday, and other Hindu holidays are also widely celebrated.
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+
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+ Muslims represent 4.97% of the population.[3] Eid al-Fitr is a public holiday and Eid al-Adha, Mawlid, Hosay, and other Muslim holidays are also celebrated. There has also been a Jewish community on the islands for many centuries, however their numbers have never been large, with a 2007 estimating putting the Jewish population at 55 individuals.[131][132]
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+ African-derived or Afrocentric religions are also practised, notably Trinidad Orisha (Yoruba) believers (0.9%) and Rastafarians (0.27%).[3] Various aspects of traditional obeah beliefs are still commonly practised on the islands.[50]
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+
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+ Two African syncretic faiths, the Shouter or Spiritual Baptists and the Orisha faith (formerly called Shangos, a less than complimentary term)[citation needed] are among the fastest growing religious groups. Similarly, there is a noticeable increase in numbers of Evangelical Protestant and Fundamentalist churches usually lumped as "Pentecostal" by most Trinidadians, although this designation is often inaccurate. Sikhism, Jainism, Bahá'í, and Buddhism are practised by a minority of Indo-Trinidadian and Tobagonians. Several eastern religions such as Buddhism and Chinese folk religions such as Taoism and Confucianism are followed by Chinese Trinidadian and Tobagonian.
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+
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+ Children generally start pre-school at two and a half years but this is not mandatory. They are however, expected to have basic reading and writing skills when they commence primary school. Students proceed to a primary school at the age of five years. Seven years are spent in primary school. The seven classes of primary school consists of First Year and Second Year, followed by Standard One through Standard Five. During the final year of primary school, students prepare for and sit the Secondary Entrance Assessment (SEA) which determines the secondary school the child will attend.
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+ Students attend secondary school for a minimum of five years, leading to the CSEC (Caribbean Secondary Education Certificate) examinations, which is the equivalent of the British GCSE O levels. Children with satisfactory grades may opt to continue high school for a further two-year period, leading to the Caribbean Advanced Proficiency Examinations (CAPE), the equivalent of GCE A levels. Both CSEC and CAPE examinations are held by the Caribbean Examinations Council (CXC). Public Primary and Secondary education is free for all, although private and religious schooling is available for a fee.
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+
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+ Tertiary education for tuition costs are provided for via GATE (The Government Assistance for Tuition Expenses), up to the level of the bachelor's degree, at the University of the West Indies (UWI), the University of Trinidad and Tobago (UTT), the University of the Southern Caribbean (USC), the College of Science, Technology and Applied Arts of Trinidad and Tobago (COSTAATT) and certain other local accredited institutions. Government also currently subsidises some Masters programmes. Both the Government and the private sector also provide financial assistance in the form of academic scholarships to gifted or needy students for study at local, regional or international universities.
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+ While women account for only 49% of the population, they constitute nearly 55% of the workforce in the country.[137]
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+ Trinidad and Tobago is the most developed nation and one of the wealthiest in the Caribbean and is listed in the top 40 (2010 information) of the 70 high-income countries in the world.[citation needed] Its gross national income per capita of US$20,070[138] (2014 gross national income at Atlas Method) is one of the highest in the Caribbean.[139] In November 2011, the OECD removed Trinidad and Tobago from its list of developing countries.[140] Trinidad's economy is strongly influenced by the petroleum industry. Tourism and manufacturing are also important to the local economy. Tourism is a growing sector, particular on Tobago, although proportionately it is much less important than in many other Caribbean islands. Agricultural products include citrus and cocoa. It also supplies manufactured goods, notably food, beverages, and cement, to the Caribbean region.
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+ Trinidad and Tobago is the leading Caribbean producer of oil and gas, and its economy is heavily dependent upon these resources.[16] Oil and gas account for about 40% of GDP and 80% of exports, but only 5% of employment.[15] Recent growth has been fuelled by investments in liquefied natural gas (LNG), petrochemicals, and steel. Additional petrochemical, aluminium, and plastics projects are in various stages of planning.
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+
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+ The country is also a regional financial centre, and the economy has a growing trade surplus.[90] The expansion of Atlantic LNG over the past six years created the largest single-sustained phase of economic growth in Trinidad and Tobago. The nation is an exporter of LNG and supplied a total of 13.4 billion m3 in 2017. The largest markets for Trinidad and Tobago's LNG exports are Chile and the United States.[141]
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+
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+ Trinidad and Tobago has transitioned from an oil-based economy to a natural gas based economy. In 2017, natural gas production totalled 18.5 billion m3, a decrease of 0.4% from 2016 with 18.6 billion m3 of production.[141] Oil production has decreased over the past decade from 7.1 million metric tonnes per year in 2007 to 4.4 million metric tonnes per year in 2017.[142] In December 2005, the Atlantic LNG's fourth production module or "train" for liquefied natural gas (LNG) began production. Train four has increased Atlantic LNG's overall output capacity by almost 50% and is the largest LNG train in the world at 5.2 million tons/year of LNG.[citation needed]
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+ Trinidad and Tobago is far less dependent on tourism than many other Caribbean countries and territories, with the bulk of tourist activity occurring on Tobago.[16] The government has made efforts to boost this sector in recent years.[16]
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+
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+ Historically agricultural production (for example, sugar and coffee) dominated the economy, however this sector has been in steep decline since the 20th century and now forms just 0.4% of the country's GDP, employing 3.1% of the workforce.[15][16] Various fruits and vegetables are grown, such as cucumbers, eggplant, cassava, pumpkin, dasheen (taro) and coconut; fishing is still also commonly practised.[15]
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+
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+ Trinidad and Tobago, in an effort to undergo economic transformation through diversification,[15] formed InvesTT in 2012 to serve as the country's sole investment promotion agency. This agency is aligned to the Ministry of Trade and Industry and is to be the key agent in growing the country's non-oil and gas sectors significantly and sustainably.[143]
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+
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+ Trinidad and Tobago has a well developed communications sector. The telecommunications and broadcasting sectors generated an estimated TT$5.63 billion (US$0.88 billion) in 2014, which as a percentage of GDP equates to 3.1 percent. This represented a 1.9 percent increase in total revenues generated by this industry compared to last year. Of total telecommunications and broadcasting revenues, mobile voice services accounted for the majority of revenues with TT$2.20 billion (39.2 percent). This was followed by internet services which contributed TT$1.18 billion or 21.1 percent. The next highest revenue earners for the industry were fixed voice services and paid television services whose contributions totalled TT$0.76 billion and TT$0.70 billion respectively (13.4 percent and 12.4 percent). International voice services was next in line, generating TT$0.27 billion (4.7 percent) in revenues. Free-to Air radio and television services contributed TT$0.18 billion and TT$0.13 billion respectively (3.2 percent and 2.4 percent). Finally, other contributors included "other revenues" and "leased line services" with earnings of TT$0.16 billion and TT$0.05 billion respectively, with 2.8 percent and 0.9 percent.[144]
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+
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+ There are several providers for each segment of the telecommunications market. Fixed Lines Telephone service is provided by Digicel, TSTT (operating as bmobile) and Cable & Wireless Communications operating as FLOW; cellular service is provided by TSTT (operating as bmobile) and Digicel whilst internet service is provided by TSTT, FLOW, Digicel, Green Dot and Lisa Communications.
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+
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+ The Government of Trinidad and Tobago has recognised the creative industries as a pathway to economic growth and development. It is one of the newest, most dynamic sectors where creativity, knowledge and intangibles serve as the basic productive resource. In 2015, the Trinidad and Tobago Creative Industries Company Limited (CreativeTT) was established as a state agency under the Ministry of Trade and Industry with a mandate to stimulate and facilitate the business development and export activities of the Creative Industries in Trinidad and Tobago to generate national wealth, and, as such, the company is responsible for the strategic and business development of the three (3) niche areas and sub sectors currently under its purview – Music, Film and Fashion. MusicTT, FilmTT and FashionTT are the subsidiaries established to fulfil this mandate.
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+
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+ The transport system in Trinidad and Tobago consists of a dense network of highways and roads across both major islands, ferries connecting Port of Spain with Scarborough and San Fernando, and international airports on both islands.[16] The Uriah Butler Highway, Churchill Roosevelt Highway and the Sir Solomon Hochoy Highway links the island of Trinidad together, whereas the Claude Noel Highway is the only major highway in Tobago. Public transportation options on land are public buses, private taxis and minibuses. By sea, the options are inter-island ferries and inter-city water taxis.[145]
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+
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+ The island of Trinidad is served by Piarco International Airport located in Piarco, which opened on 8 January 1931.[citation needed] Elevated at 17.4 metres (57 ft) above sea level it comprises an area of 680 hectares (1,700 acres) and has a runway of 3,200 metres (10,500 ft). The airport consists of two terminals, the North Terminal and the South Terminal. The older South Terminal underwent renovations in 2009 for use as a VIP entrance point during the 5th Summit of the Americas. The North Terminal was completed in 2001, and consists of[146] 14-second-level aircraft gates with jetways for international flights, two ground-level domestic gates and 82 ticket counter positions.
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+
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+ In 2008 the passenger throughput at Piarco International Airport was approximately 2.6 million. It is the seventh busiest airport in the Caribbean and the third busiest in the English-speaking Caribbean, after Sangster International Airport and Lynden Pindling International Airport.[citation needed] Caribbean Airlines, the national airline, operates its main hub at the Piarco International Airport and services the Caribbean, the United States, Canada and South America. The airline is wholly owned by the Government of Trinidad and Tobago. After an additional cash injection of US$50 million, the Trinidad and Tobago government acquired the Jamaican airline Air Jamaica on 1 May 2010, with a 6–12-month transition period to follow.[147]
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+
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+ The Island of Tobago is served by the A.N.R. Robinson International Airport in Crown Point.[16] This airport has regular services to North America and Europe. There are regular flights between the two islands, with fares being heavily subsidised by the Government.
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+
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+ Trinidad was formerly home to a railway network, however this was closed down in 1968.[148] There have been talks to build a new railway on the islands, though nothing yet has come of this.[149]
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+
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+ The Strategic Plan for the Caribbean Community 2015–2019 was adopted by Trinidad and Tobago and the other members of the Caribbean Common Market (CARICOM) in 2014. The first of its kind, this document reflects a desire among countries to embrace a more profound regionalism, in order to reposition the Caribbean in an increasingly volatile global economy. The plan proposes mobilising funding from the public and private sectors to foster research and development (R&D) in the business sector. The plan outlines strategies for nurturing creativity, entrepreneurship, digital literacy and for making optimum use of available resources. It focuses on developing creative, manufacturing and service industries, with a special focus on tourism initially, natural resources and value-added products, agriculture and fisheries, to reduce dependence on food imports and foster sustainable fisheries, and energy efficiency.[150][151]
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+
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+ Trinidad and Tobago is the region's leading exporter of oil and gas but imports of fossil fuels provided over 90% of the energy consumed by its CARICOM neighbours in 2008. This vulnerability led CARICOM to develop an Energy Policy which was approved in 2013. This policy is accompanied by the CARICOM Sustainable Energy Roadmap and Strategy (C-SERMS). Under the policy, renewable energy sources are to contribute 20% of the total electricity generation mix in member states by 2017, 28% by 2022 and 47% by 2027.[150]
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+
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+ In 2014 Trinidad and Tobago was the third country in the world which emitted the most CO2 per capita after Qatar and Curacao according to the World Bank.[152] On average, each inhabitant produced 34.2 metric tons of CO2 in the atmosphere. In comparison, the world average was 5.0 tons per capita the same year.
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+
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+ The Caribbean Community Climate Change Centre (CCCCC) has produced an implementation plan for 2011–2021 and carried out work to assess and build capacity in climate change mitigation and resilient development strategies. This work has been supported by the region's specialists, who have produced models for climate change and mitigation processes in Caribbean states. They also play a major advisor role to the divisions in ministries responsible for climate change. The growing frequency and intensity of hurricanes is of concern to all Caribbean nations. In 2012, Trinidad and Tobago had a 9% chance each year of being struck by a hurricane, according to estimates by the International Monetary Fund.[150][153]
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+ The two main bodies responsible for science, technology and innovation in Trinidad and Tobago are the Ministry of Science, Technology and Higher Education and the National Commission for Science and Technology.[150]
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+
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+ In 2002, Trinidad and Tobago adopted Vision 2020. Like Jamaica's Vision 2030 (2009) and the Strategic Plan of Barbados for 2005–2025, Trinidad and Tobago's Vision 2020 accords central importance to harnessing science, technology and innovation (STI) to raise living standards and strengthen resilience to environmental shocks like hurricanes.[150]
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+
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+ Trinidad and Tobago is one of the more affluent members of CARICOM, thanks to its natural resources. Despite this, it spent just 0.05% of GDP on R&D in 2012, according to the UNESCO Institute for Statistics. Even when the country was enjoying economic growth of 8% per annum in 2004, it devoted just 0.11% of GDP to R&D. Calculated in thousands of current Purchasing Power Parity (PPP) dollars, research expenditure actually dropped between 2009 and 2012 from 21 309 to 19 232. This corresponds to research expenditure of $PPP 65 per capita in 2009 and $PPP 45 in 2012.[150]
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+
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+ Industrial R&D has declined since 2000, perhaps owing to the drop in research activity in the sugar sector. Whereas industrial R&D accounted for 24% of domestic research expenditure in 2004 and 29.5% in 2005, it had become almost non-existent by 2010.[150]
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+
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+ The number of researchers in Trinidad and Tobago grew from 787 to 914 between 2009 and 2012. This corresponds in a rise from 595 to 683 in the number of researchers (head counts) per million inhabitants.[150]
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+
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+ Scientific output grew between 2007 and 2011, according to Thomson Reuters' Web of Science (Science Citation Index Expanded) before contracting over the period 2012–2014. Trinidad and Tobago produced 109 publications per million population in 2014, behind Grenada (1,430), St Kitts and Nevis (730), Barbados (182) and Dominica (138) but ahead of the Bahamas (86), Belize (47) and Jamaica (42).[150]
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+
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+ Between 2008 and 2014, scientists collaborated most with their peers from the United States (251 papers), United Kingdom (183), Canada (95), India (63) and Jamaica (43), according to the copublication record of Thomson Reuters. In turn, Jamaican scientists considered their counterparts from Trinidad and Tobago to be their fourth-closest collaborators (with 43 joint papers) after those from the United States, United Kingdom and Canada.[150]
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+
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+ Between 2008 and 2013, Trinidad and Tobago registered 17 patents with the US Patent and Trademark Office (USPTO). This corresponds to 13% of the 134 patents registered by CARICOM members over this period. The top contributors were the Bahamas (34 patents) and Jamaica (22).[150]
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+ Trinidad and Tobago led CARICOM members for the value of high-tech exports in 2008 (US$36.2 million) but these exports plummeted to US$3.5 million the following year, according to the Comtrade database of the United Nations Statistics Division.[150]
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+
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+ The Caribbean Industrial Research Institute in Trinidad and Tobago facilitates climate change research and provides industrial support for R&D related to food security. It also carries out equipment testing and calibration for major industries.[150]
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+ The Caribbean Epidemiology Centre in Port of Spain, University of Trinidad and Tobago, Tobago Institute of Health, and University of the West Indies (St Augustine campus) also conduct R&D.[150]
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+ Probability of a hurricane striking Caribbean countries in a given year, 2012 (%).[154]
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+ Electricity costs for the CARICOM countries, 2011[155]
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+ GERD by sector of performance in Trinidad and Tobago, 2000–2012[156]
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+ Scientific publication trends in the CARICOM countries, 2005–2014[157]
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+ Scientific publications in the CARICOM countries, 2014[157]
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+ USPTO patents granted to Caribbean countries, 2008–2013.[158]
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+ Trinidad and Tobago has a diverse culture mixing Indian, African, Creole, Chinese, Amerindian, Arab, Latino, and European influences, reflecting the various communities who have migrated to the islands over the centuries. The island is particularly renowned for its annual Carnival celebrations.[16] Festivals rooted in various religions and cultures practiced on the islands are also popular, such as Christmas, Divali, Phagwah (Holi), Easter, New Year’s Day, Hosay, Eid al-Fitr, the Santa Rosa Indigenous Festival, and Chinese New Year.
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+ Trinidad and Tobago claims two Nobel Prize-winning authors, V. S. Naipaul and St Lucian-born Derek Walcott (who also founded the Trinidad Theatre Workshop). Other notable writers include Neil Bissoondath, Vahni Capildeo, Earl Lovelace, Seepersad Naipaul, Shiva Naipaul, Lakshmi Persaud, Kenneth Ramchand, Arnold Rampersad, and Samuel Selvon.
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+ Trinidadian designer Peter Minshall is renowned not only for his Carnival costumes but also for his role in opening ceremonies of the Barcelona Olympics, the 1994 FIFA World Cup, the 1996 Summer Olympics, and the 2002 Winter Olympics, for which he won an Emmy Award.[161]
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+ Trinidad and Tobago is the birthplace of calypso music and the steelpan.[162][163][164] Trinidad is also the birthplace of soca music, chutney music, chutney-soca, parang, rapso, pichakaree and chutney parang.
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+ The limbo dance originated in Trinidad as an event that took place at wakes in Trinidad. The limbo has African roots. It was popularized in the 1950s by dance pioneer Julia Edwards[165] (known as the First Lady of Limbo) and her company which appeared in several films.[166] Bélé, Bongo, and whining are also dance forms with African roots.[167]
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+ Jazz, ballroom, ballet, modern, and salsa dancing are also popular.[167]
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+ Indian dance forms are also popular in Trinidad and Tobago.[168] Kathak, Odissi, and Bharatanatyam are the most popular Indian classical dance forms in Trinidad and Tobago.[169] Indian folk dances and Bollywood dances are also popular.[169]
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+ Geoffrey Holder (brother of Boscoe Holder) and Heather Headley are two Trinidad-born artists who have won Tony Awards for theatre. Holder also has a distinguished film career, and Headley has won a Grammy Award as well.
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+ Indian theatre is also popular throughout Trinidad and Tobago. Dramas such as Nautankis, Raja Harishchandra, Sharwan Kumar, and Alha-Khand were brought by Indians to Trinidad and Tobago, however they had largely began to die out, till preservation began by Indian cultural groups.[170] The drama about the life of the Hindu god Rama, Ramleela, is popular during the time between Sharad Navaratri and Dushera and the drama about the life of the Hindu god Krishna, Ras leela (Krishna leela), is popular around the time of Krishna Janmashtami.[171][172][173]
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+ Trinidad and Tobago is also smallest country to have two Miss Universe titleholders and the first black woman ever to win: Janelle Commissiong in 1977, followed by Wendy Fitzwilliam in 1998; the country has also had one Miss World titleholder, Giselle LaRonde.[citation needed]
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+ Hasely Crawford won the first Olympic gold medal for Trinidad and Tobago in the men's 100-metre dash in the 1976 Summer Olympics. Nine different athletes from Trinidad and Tobago have won twelve medals at the Olympics, beginning with a silver medal in weightlifting, won by Rodney Wilkes in 1948,[174] and most recently, a gold medal by Keshorn Walcott in the men's javelin throw in 2012. Ato Boldon has won the most Olympic and World Championship medals for Trinidad and Tobago in athletics, with eight in total – four from the Olympics and four from the World Championships. Boldon was the sole world champion Trinidad and Tobago has produced until Jehue Gordon in Moscow 2013. Ato won the 1997 200 m sprint World Championship in Athens. Swimmer George Bovell III won a bronze medal in the men's 200 m IM in 2004. At the 2017 World Championship in London, the Men 4x400 relay team captured the title, thus the country now celebrates three world championships titles. The team consisted of Jarrin Solomon, Jareem Richards, Machel Cedenio and Lalonde Gordon with Renny Quow who ran in the heats.
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+ Also in 2012 Lalonde Gordon competed in the XXX Summer Olympics where he won a bronze medal in the 400 metres (1,300 feet), being surpassed by Luguelin Santos of the Dominican Republic and Kirani James of Grenada. Keshorn Walcott (as stated above) came first in javelin and earned a gold medal, making him the second Trinidadian in the country's history to receive one. This also makes him the first Western[clarification needed] athlete in 40 years to receive a gold medal in the javelin sport, and the first athlete from Trinidad and Tobago to win a gold medal in a field event in the Olympics. Sprinter Richard Thompson is also from Trinidad and Tobago. He came second place to Usain Bolt in the Beijing Olympics in the 100 metres (330 feet) with a time of 9.89s.
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+ In 2018 The Court of Arbitration for Sport made its final decision on the failed doping sample from the Jamaican team in the 4 x 100 relay in the 2008 Olympic Games. The team from Trinidad and Tobago will be awarded the gold medal, because of the second rank during the relay run.[175]
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+ Cricket is a popular sport of Trinidad and Tobago, often deemed the national sport, and there is intense inter-island rivalry with its Caribbean neighbours. Trinidad and Tobago is represented at Test cricket, One Day International as well as Twenty20 cricket level as a member of the West Indies team. The national team plays at the first-class level in regional competitions such as the Regional Four Day Competition and Regional Super50. Meanwhile, the Trinbago Knight Riders play in the Caribbean Premier League.
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+ The Queen's Park Oval located in Port of Spain is the largest cricket ground in the West Indies, having hosted 60 Test matches as of January 2018. Trinidad and Tobago along with other islands from the Caribbean co-hosted the 2007 Cricket World Cup.
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+ Brian Lara, world record holder for the most runs scored both in a Test and in a First Class innings and other records, was born in a small town of Santa Cruz and is often referred to as the Prince of Port of Spain or simply the Prince. This legendary West Indian batsman is widely regarded (along with Sir Donald Bradman, Sunil Gavaskar and Sachin Tendulkar[citation needed]) as one of the best batsmen ever to have played the game,[citation needed] and is one of the most famous sporting icons in the country.
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+ Association football is also a popular sport in Trinidad and Tobago. The men's national football team qualified for the 2006 FIFA World Cup for the first time by beating Bahrain in Manama on 16 November 2005, making them the second smallest country ever (in terms of population) to qualify, after Iceland. The team, coached by Dutchman Leo Beenhakker, and led by Tobagonian-born captain Dwight Yorke, drew their first group game – against Sweden in Dortmund, 0–0, but lost the second game to England on late goals, 0–2. They were eliminated after losing 2–0 to Paraguay in the last game of the Group stage. Prior to the 2006 World Cup qualification, Trinidad and Tobago came close in a controversial qualification campaign for the 1974 FIFA World Cup. Following the match, the referee of their critical game against Haiti was awarded a lifetime ban for his actions.[176] Trinidad and Tobago again fell just short of qualifying for the World Cup in 1990, needing only a draw at home against the United States but losing 1–0.[177] They play their home matches at the Hasely Crawford Stadium. Trinidad and Tobago hosted the 2001 FIFA U-17 World Championship, and hosted the 2010 FIFA U-17 Women's World Cup.
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+ The TT Pro League is the country's primary football competition and is the top level of the Trinidad and Tobago football league system. The Pro League serves as a league for professional football clubs in Trinidad and Tobago. The league began in 1999 as part of a need for a professional league to strengthen the country's national team and improve the development of domestic players. The first season took place in the same year beginning with eight teams.
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+ Basketball is commonly played in Trinidad and Tobago in colleges, universities and throughout various urban basketball courts. Its national team is one of the most successful teams in the Caribbean. At the Caribbean Basketball Championship it won four straight gold medals from 1986 to 1990.
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+ Netball has long been a popular sport in Trinidad and Tobago, although it has declined in popularity in recent years. At the Netball World Championships they co-won the event in 1979, were runners up in 1987, and second runners up in 1983.
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+ Rugby is played in Trinidad and Tobago and continues to be a popular sport, and horse racing is regularly followed in the country.
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+ There is also the Trinidad and Tobago national baseball team which is controlled by the Baseball/Softball Association of Trinidad and Tobago, and represents the nation in international competitions. The team is a provisional member of the Pan American Baseball Confederation.
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+ There are a number of 9 and 18-hole golf courses on Trinidad and Tobago. The most established is the St Andrews Golf Club, Maraval in Trinidad (commonly referred to as Moka), and there is a newer course at Trincity, near Piarco Airport called Millennium Lakes. There are 18-hole courses at Chaguramas and Point-a-Pierre and 9-hole courses at Couva and St Madeline. Tobago has two 18-hole courses. The older of the two is at Mount Irvine, with the Magdalena Hotel & Golf Club (formerly Tobago Plantations) being built more recently.
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+ Although a minor sport, bodybuilding is of growing interest in Trinidad and Tobago. Heavyweight female bodybuilder Kashma Maharaj is of Trinidadian descent.
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+ Dragonboat is also another water-sport that has been rapidly growing over the years. Introduced in 2006. the fraternity made consistent strides in having more members apart of the TTDBF (Trindad and Tobago Dragonboat Federation) as well as performing on an international level such as the 10th IDBF World Nations Dragon Boat Championships in Tampa, Florida in the US in 2011.
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+ Claude Noel is a former world champion in professional boxing. He was born in Tobago.
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+ The flag was chosen by the Independence committee in 1962. Red, black and white symbolise the warmth of the people, the richness of the earth and water respectively.[178][179]
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+ The coat of arms was designed by the Independence committee, and features the scarlet ibis (native to Trinidad), the cocrico (native to Tobago) and hummingbird. The shield bears three ships, representing both the Trinity, and the three ships that Columbus sailed.[178]
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+ There are five categories and thirteen classes of national awards:[180]
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+ The national anthem of the twin-island state is "Forged from the Love of Liberty".[181][182]
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+ Other national songs include "God Bless Our Nation"[183] and "Our Nation's Dawning".[184]
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+ The national flower of Trinidad and Tobago is the chaconia flower. It was chosen as the national flower because it is an indigenous flower that has witnessed the history of Trinidad and Tobago. It was also chosen as the national flower because of its red colour that resembles the red of the national flag and coat of arms and because it blooms around the Independence Day of Trinidad and Tobago.[185]
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+ The national birds of Trinidad and Tobago are the scarlet ibis and the cocrico. The scarlet ibis is kept safe by the government by living in the Caroni Bird Sanctuary which was set up by the government for the protection of these birds. The Cocrico is more indigenous to the island of Tobago and are more likely to be seen in the forest.[186] The hummingbird is considered another symbol of Trinidad and Tobago due to its significance to the indigenous peoples, however, it is not a national bird.[187]
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+ The scarlet ibis birds flying over the Caroni Swamp
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+ The cocrico bird in Tobago
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+ This article incorporates text from a free content work. UNESCO Science Report: towards 2030, 156–173, Harold Ramkissoon & Ishenkumba A. Kahwa, UNESCO Publishing. To learn how to add open license text to Wikipedia articles, please see this how-to page. For information on reusing text from Wikipedia, please see the terms of use.
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+ ^ These three Dutch Caribbean territories form the SSS islands.
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+ * These three Dutch Caribbean territories form the BES islands.
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+ † Physiographically, these are continental islands not a part of the volcanic Windward Islands arc. However, based on proximity, these islands are sometimes grouped with the Windward Islands culturally and politically.
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+ ~ Disputed territories administered by Colombia.
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+ # Physiographically, Bermuda is an isolated oceanic island in the North Atlantic Ocean, not a part of the Antilles, West Indies, Caribbean, North American continent or South American continent. Usually grouped with Northern American countries based on proximity; occasionally grouped with the Caribbean region culturally.
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+ Coordinates: 10°36′N 61°6′W / 10.600°N 61.100°W / 10.600; -61.100
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+ Trinidad and Tobago (/ˈtrɪnɪdæd ... təˈbeɪɡoʊ/ (listen), /- toʊ-/), officially the Republic of Trinidad and Tobago, is the southernmost island country in the Caribbean.[14][15] Consisting of the main islands Trinidad and Tobago, and numerous much smaller islands, it is situated 130 kilometres (81 miles) south of Grenada and 11 kilometres (6.8 miles) off the coast of northeastern Venezuela.[16] It shares maritime boundaries with Barbados to the northeast, Grenada to the northwest, Guyana to the southeast, and Venezuela to the south and west.[17][18]
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+ The island of Trinidad was inhabited for centuries by native Amerindian peoples before becoming a colony in the Spanish Empire, following the arrival of Christopher Columbus in 1498. Spanish governor Don José María Chacón surrendered the island to a British fleet under the command of Sir Ralph Abercromby in 1797.[19] During the same period, the island of Tobago changed hands among Spanish, British, French, Dutch and Courlander colonisers more times than any other island in the Caribbean.[citation needed] Trinidad and Tobago were ceded to Britain in 1802 under the Treaty of Amiens as separate states and unified in 1889.[20] Trinidad and Tobago obtained independence in 1962, becoming a republic in 1976.[15][16]
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+ Trinidad and Tobago has the third highest GDP per capita based on purchasing power parity (PPP) in the Americas after the United States and Canada.[21] It is recognised by the World Bank as a high-income economy.[22] Unlike most Caribbean nations and territories, which rely heavily on tourism, the Trinidadian economy is primarily industrial with an emphasis on petroleum and petrochemicals;[23] much of the nation's wealth is derived from its large reserves of oil and natural gas.[24][25]
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+ Trinidad and Tobago is well known for its African and Indian cultures, reflected in its large and famous Carnival, Diwali, and Hosay celebrations, as well being the birthplace of steelpan, the limbo, and music styles such as calypso, soca, rapso, parang, chutney, and chutney soca.[26][27][28][29][30][31][32][33]
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+ Historian E. L. Joseph claimed that Trinidad's Amerindian name was Cairi or "Land of the Humming Bird", derived from the Arawak name for hummingbird, ierèttê or yerettê.[citation needed] However, other authors dispute this etymology with some claiming that cairi does not mean hummingbird (tukusi or tucuchi being suggested as the correct word) and some claiming that kairi, or iere, simply means island.[clarification needed][34][citation needed] Christopher Columbus renamed it "La Isla de la Trinidad" ("The Island of the Trinity"), fulfilling a vow made before setting out on his third voyage of exploration.[35] Tobago's cigar-like shape, or the use of tobacco by the native people, may have given it its Spanish name (cabaco, tavaco, tobacco) and possibly some of its other Amerindian names, such as Aloubaéra (black conch) and Urupaina (big snail),[34] although the English pronunciation is /təˈbeɪɡoʊ/.[citation needed]
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+ Both Trinidad and Tobago were originally settled by Amerindians who came through South America.[16] Trinidad was first settled by pre-agricultural Archaic people at least 7,000 years ago, making it the earliest settled part of the Caribbean.[36] Banwari Trace in south-west Trinidad is the oldest attested archaeological site in the Caribbean, dating to about 5000 BC. Several waves of migration occurred over the following centuries, which can be identified by differences in their archaeological remains.[37] At the time of European contact, Trinidad was occupied by various Arawakan-speaking groups including the Nepoya and Suppoya, and Cariban-speaking groups such as the Yao, while Tobago was occupied by the Island Caribs and Galibi. Trinidad was known to the native peoples as 'Ieri' ('Land of the Humming Bird').[36]
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+ Christopher Columbus was the first European to see Trinidad, on his third voyage to the Americas in 1498.[36][38] He also reported seeing Tobago on the distant horizon, naming it Bellaforma, but did not land on the island.[16][39]
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+ In the 1530s Antonio de Sedeño, a Spanish soldier intent on conquering the island of Trinidad, landed on its southwest coast with a small army of men, intending to subdue the Amerindian peoples of the island. Sedeño and his men fought the native peoples on many occasions, and subsequently built a fort. The next few decades were generally spent in warfare with the native peoples, until in 1592, the 'Cacique' (native chief) Wannawanare (also known as Guanaguanare) granted the area around modern Saint Joseph to Domingo de Vera e Ibargüen, and withdrew to another part of the island.[34] The settlement of San José de Oruña was later established by Antonio de Berrío on this land in 1592.[16][36] Shortly thereafter the English sailor Sir Walter Raleigh arrived in Trinidad on 22 March 1595 in search of the long-rumoured "El Dorado" ('City of Gold') supposedly located in South America.[36] He attacked San José, captured and interrogated Antonio de Berrío, and obtained much information from him and from the Cacique Topiawari; Raleigh then went on his way, and Spanish authority was restored.[40]
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+ Meanwhile, there were numerous attempts by European powers to settle Tobago during the 1620-40s, with the Dutch, English and Courlanders (people from the Duchy of Courland and Semigallia, now part of Latvia) all attempting to colonise the island with little success.[41][42] From 1654 the Dutch and Courlanders managed to gain a more secure foothold, later joined by several hundred French settlers.[41] A plantation economy developed based on the production of sugar, indigo and rum, worked by large numbers of African slaves who soon came to vastly outnumber the European colonists.[42][41] Large numbers of forts were constructed as Tobago became a source of contention between France, Holland and Britain, with the island changing hands some 31 times prior to 1814, a situation exacerbated by widespread piracy.[42] The British managed to hold Tobago from 1762–1781, whereupon it was captured by the French, who ruled until 1793 when Britain re-captured the island.[42]
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+ The 17th century on Trinidad passed largely without major incident, but sustained attempts by the Spaniards to control and rule over the Amerindians were often fiercely resisted.[36] In 1687 the Catholic Catalan Capuchin friars were given responsibility for the conversions of the indigenous people of Trinidad and the Guianas.[36] They founded several missions in Trinidad, supported and richly funded by the state, which also granted encomienda right to them over the native peoples, in which the native peoples were forced to provide labour for the Spanish.[36] One such mission was Santa Rosa de Arima, established in 1789, when Amerindians from the former encomiendas of Tacarigua and Arauca (Arouca) were relocated further west.[citation needed] Escalating tensions between the Spaniards and Amerindians culminated in violence 1689, when Amerindians in the San Rafael encomienda rebelled and killed several priests, attacked a church, and killed the Spanish governor José de León y Echales. Among those killed in the governor's party was Juan Mazien de Sotomayor, missionary priest to the Nepuyo villages of Cuara, Tacarigua and Arauca.[citation needed] The Spanish retaliated severely, slaughtering hundreds of native peoples in an event that became known as the Arena massacre.[36] As a result of this, continuing Spanish slave-raiding, and the devastating impact of introduced disease to which they had no immunity, the native population was virtually wiped out by the end of the following century.[43][36]
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+ During this period Trinidad was an island province belonging to the Viceroyalty of New Spain, together with Central America, present-day Mexico and the southwestern United States.[44] In 1757 the capital was moved from San José de Oruña to Puerto de España (modern Port of Spain) following several pirate attacks.[45] However the Spanish never made any concerted effort to colonise the islands; Trinidad in this period was still mostly forest, populated by a few Spaniards with a handful of slaves and a few thousand Amerindians.[44] Indeed, the population in 1777 was only 1,400, and Spanish colonisation in Trinidad remained tenuous.[citation needed]
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+ Since Trinidad was considered underpopulated, Roume de St. Laurent, a Frenchman living in Grenada, was able to obtain a Cédula de Población from the Spanish king Charles III on 4 November 1783.[46] A Cédula de Población had previously been granted in 1776 by the king, but had not shown results, and therefore the new Cédula was more generous.[16] It granted free land and tax exemption for 10 years to Roman Catholic foreign settlers who were willing to swear allegiance to the King of Spain.[16] The Spanish also gave many incentives to lure settlers to the island, including exemption from taxes for ten years and land grants in accordance with the terms set out in the Cédula.[47] The land grant was 30 fanegas (13 hectares/32 acres) for each free man, woman and child and half of that for each slave that they brought with them. The Spanish sent a new governor, José María Chacón, to implement the terms of the new cédula.[46]
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+ It was fortuitous that the Cédula was issued only a few years before the French Revolution. During that period of upheaval, French planters with their slaves, free coloureds and mulattos from the neighbouring islands of Martinique, Saint Lucia, Grenada, Guadeloupe and Dominica migrated to Trinidad, where they established an agriculture-based economy (sugar and cocoa).[44] These new immigrants established local communities in Blanchisseuse, Champs Fleurs, Paramin,[48] Cascade, Carenage and Laventille.
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+
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+ As a result, Trinidad's population jumped to over 15,000 by the end of 1789, and by 1797 the population of Port of Spain had increased from under 3,000 to 10,422 in just five years, with a varied population of mixed race individuals, Spaniards, Africans, French republican soldiers, retired pirates and French nobility.[44] The total population of Trinidad was 17,718, of which 2,151 were of European ancestry, 4,476 were "free blacks and people of colour", 10,009 were enslaved people and 1,082 Amerindians.[citation needed] The sparse settlement and slow rate of population-increase during Spanish rule (and even later during British rule) made Trinidad one of the less populated colonies of the West Indies, with the least developed plantation infrastructure.[49]
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+
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+ The British had begun to take a keen interest in Trinidad, and in 1797 a British force led by General Sir Ralph Abercromby launched an invasion of Trinidad.[16][50] His squadron sailed through the Bocas and anchored off the coast of Chaguaramas. Seriously outnumbered, Governor Chacón decided to capitulate to British without fighting.[50] Trinidad thus became a British crown colony, with a largely French-speaking population and Spanish laws.[44] British rule was later formalised under the Treaty of Amiens (1802).[16][50] The colony's first British governor was Thomas Picton, however his heavy-handed approach to enforcing British authority, including the use of torture and arbitrary arrest, led to his being recalled.[50]
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+ British rule led to an influx of settlers from the United Kingdom and the British colonies of the Eastern Caribbean. English, Scots, Irish, German and Italian families arrived, as well as some free blacks known as 'Merikins' who had fought for Britain in the War of 1812 and were granted land in southern Trinidad.[51][52][53] Under British rule, new states were created and the importation of slaves increased, however by this time support for abolitionism had vastly increased and in England the slave trade was under attack.[49][54] Slavery was abolished in 1833, after which former slaves served an "apprenticeship" period. In 1837 Daaga, a West African slave trader who had been captured by Portuguese slavers and later rescued by the British navy, was conscripted into the local regiment. Daaga and a group of his compatriots mutinied at the barracks in St Joseph and set out eastward in an attempt to return to their homeland. The mutineers were ambushed by a militia unit just outside the town of Arima. The revolt was crushed at the cost of some 40 dead, and Daaga and his party were later executed at St Joseph.[55] The apprenticeship system ended on 1 August 1838 with full emancipation.[16][53] An overview of the populations statistics in 1838, however, clearly reveals the contrast between Trinidad and its neighbouring islands: upon emancipation of the slaves in 1838, Trinidad had only 17,439 slaves, with 80% of slave owners having enslaved fewer than 10 people each.[56] In contrast, at twice the size of Trinidad, Jamaica had roughly 360,000 slaves.[57]
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+
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+ After the African slaves were emancipated many refused to continue working on the plantations, often moving out to urban areas such as Laventille and Belmont to the east of Port of Spain.[53] As a result, a severe agricultural labour shortage emerged; the British filled this gap by instituting a system of indentureship. Various nationalities were contracted under this system, including Indians, Chinese, and Portuguese.[58] Of these, the East Indians were imported in the largest numbers, starting from 1 May 1845, when 225 Indians were brought in the first shipment to Trinidad on the Fatel Razack, a Muslim-owned vessel.[53][59] Indentureship of the Indians lasted from 1845 to 1917, during which time more than 147,000 Indians came to Trinidad to work on sugarcane plantations.[16][60]
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+
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+ Indentureship contracts were sometimes exploitative, to such an extent that historians such as Hugh Tinker were to call it "a new system of slavery". Despite these descriptions, it was not truly a new form of slavery, as workers were paid, contracts were finite, and the idea of an individual being another's property had been eliminated when slavery was abolished.[61] In addition, employers of indentured labour had no legal right to flog or whip their workers; the main legal sanction for the enforcement of the indenture laws was prosecution in the courts, followed by fines or (more likely) jail sentences.[62] People were contracted for a period of five years, with a daily wage as low as 25 cents in the early 20th century, and they were guaranteed return passage to India at the end of their contract period. However, coercive means were often used to retain labourers, and the indentureship contracts were soon extended to 10 years from 1854 after the planters complained that they were losing their labour too early.[49][53] In lieu of the return passage, the British authorities soon began offering portions of land to encourage settlement, and by 1902, more than half of the sugar cane in Trinidad was being produced by independent cane farmers; the majority of which were Indians.[63] Despite the trying conditions experienced under the indenture system, about 90% of the Indian immigrants chose, at the end of their contracted periods of indenture, to make Trinidad their permanent home.[63] East Indians entering the colony were also subject to certain crown laws which segregated them from the rest of Trinidad's population, such as the requirement that they carry a pass with them if they left the plantations, and that if freed, they carry their "Free Papers" or certificate indicating completion of the indenture period.[64]
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+ Few Indians settled on Tobago however, and the descendants of African slaves continued to form the majority of the island's population. An ongoing economic slump in the middle-to-late 19th century caused
40
+ widespread poverty.[65] Discontent erupted into rioting on the Roxborough plantation in 1876, in an event known as the Belmanna Uprising after a policeman who was killed.[65] The British eventually managed to restore control, however as a result of the disturbances Tobago's Legislative Assembly voted to dissolve itself and the island became a Crown colony in 1877.[65] With the sugar industry in a state of near-collapse and the island no longer profitable, the British attached Tobago to their Trinidad colony in 1899.[16][66][67]
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+ In 1903, a protest against the introduction of new water rates in Port of Spain erupted into rioting; 18 people were shot dead, and the Red House (the government headquarters) was damaged by fire.[66] A local elected assembly with some limited powers was introduced in 1913.[66] Economically Trinidad and Tobago remained a predominantly agricultural colony; alongside sugarcane, the cacao (cocoa) crop also contributed greatly to economic earnings in the late 19th and early 20th centuries.
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+
44
+ In November 1919, the dockworkers went on strike over bad management practises, low wages compared to a higher cost of living.[68] Strikebreakers were brought in to keep a minimum of goods moving through the ports. On December 1, 1919, the striking dockworkers rushed the harbour and chased off the strikebreakers.[68] They then proceeded to march on the government buildings in Port of Spain. Other unions and workers, many with the same grievances, joined the dock worker's strike making it a General Strike.[68] Violence broke out and was only put down with help from the sailors of British Naval ship HMS Calcutta. The unity brought upon by the strike was the first time of cooperation between the various ethnic groups of the time.[69] Historian Brinsley Samaroo says that the 1919 strikes "seem to indicate that there was a growing class consciousness after the war and this transcended racial feelings at times."[69]
45
+
46
+ However, in the 1920s, the collapse of the sugarcane industry, concomitant with the failure of the cocoa industry, resulted in widespread depression among the rural and agricultural workers in Trinidad, and encouraged the rise of a labour movement. Conditions on the islands worsened in the 1930s with the onset of the Great Depression, with an outbreak of labour riots occurring in 1937 which resulted in several deaths.[70] The labour movement aimed to unite the urban working class and agricultural labour class; the key figures being Arthur Cipriani, who led the Trinidad Workingmen's Association (TWA), and Tubal Uriah "Buzz" Butler of the British Empire Citizens' and Workers' Home Rule Party.[70] As the movement developed calls for greater autonomy from British colonial rule became widespread; this effort was severely undermined by the British Home Office and by the British-educated Trinidadian elite, many of whom were descended from the plantocracy class.
47
+
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+ Petroleum had been discovered in 1857, but became economically significant only in the 1930s and afterwards as a result of the collapse of sugarcane and cocoa, and increasing industrialisation.[71]
49
+ [72][73] By the 1950s petroleum had become a staple in Trinidad's export market, and was responsible for a growing middle class among all sections of the Trinidad population. The collapse of Trinidad's major agricultural commodities, followed by the Depression, and the rise of the oil economy, led to major changes in the country's social structure.
50
+
51
+ The presence of American military bases in Chaguaramas and Cumuto in Trinidad during World War II had a profound effect on society. The Americans vastly improved the infrastructure on Trinidad and provided many locals with well-paying jobs; however the social effects of having so many young soldiers stationed on the island, as well as their often unconcealed racial prejudice, caused resentment.[66] The Americans left in 1961.[74]
52
+
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+ In the post-war period the British began a process of decolonisation across the British Empire. In 1945 universal suffrage was introduced to Trinidad and Tobago.[16][66] Political parties emerged on the island, however these were largely divided along racial lines: Afro-Trinidadians and Tobagonians primarily supported the People's National Movement (PNM), formed in 1956 by Eric Williams, with Indo-Trinidadians and Tobagonians mostly supporting the People's Democratic Party (PDP), formed in 1953 by Bhadase Sagan Maraj,[75] which later merged into the Democratic Labour Party (DLP) in 1957.[76] Britain's Caribbean colonies formed the West Indies Federation in 1958 as a vehicle for independence, however the Federation dissolved after Jamaica withdrew following a membership referendum in 1961. The government of Trinidad and Tobago subsequently chose to seek independence from the United Kingdom on its own.[77]
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+
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+ Trinidad and Tobago gained its independence from the United Kingdom on 31 August 1962.[16][73] Elizabeth II remained head of state as Queen of Trinidad and Tobago, represented locally by Governor-General Solomon Hochoy. Eric Williams of the PNM, a noted historian and intellectual widely regarded as The Father of The Nation, became the first Prime Minister, serving in that capacity uninterrupted until 1981.[16] The dominant figure in the opposition in the early independence years was Rudranath Capildeo of the DLP. The 1960s saw the rise of a Black Power movement, inspired in part by the civil rights movement in the United States. Protests and strikes became common, with events coming to head in April 1970 when police shot dead a protester named Basil Davis.[76] Fearing a breakdown of law and order, Prime Minister Williams declared a state of emergency and arrested many of the Black Power leaders. Some army leaders who were sympathetic to the Black Power movement, notably Raffique Shah and Rex Lassalle, attempted to mutiny; however, this was quashed by the Trinidad and Tobago Coast Guard.[76] Williams and the PNM retained power, largely due to divisions in the opposition.[76]
56
+
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+ In 1963 Tobago was struck by Hurricane Flora, which killed 30 people and resulted in enormous destruction across the island.[78] Partly as a result of this, tourism came to replace agriculture as the island's main income earner in the subsequent decades.[78]
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+
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+ Between the years 1972 and 1983, the country profited greatly from the rising price of oil and the discovery of vast new oil deposits in its territorial waters, resulting in an economic boom that increased living standards greatly.[16][76] In 1976 the country became a republic within the Commonwealth, though it retained the Judicial Committee of the Privy Council as its final appellate court.[16] The position of governor-general was replaced with that of President; Ellis Clarke was the first to hold this largely ceremonial role.[79] Tobago was granted limited self-rule with the creation of the Tobago House of Assembly in 1980.[65]
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+
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+ Williams died in 1981, being replaced by George Chambers who led the country until 1986. By this time a fall in the price of oil had resulted in a recession, causing rising inflation and unemployment.[80] The main opposition parties united under the banner of National Alliance for Reconstruction (NAR) and won the 1986 Trinidad and Tobago general election, with NAR leader A. N. R. Robinson becoming the new Prime Minister.[81][76] Robinson was unable to hold together the fragile NAR coalition, and social unrest was caused by his economic reforms, such as devaluing the currency and implementing an International Monetary Fund Structural Adjustment Program.[16] In 1990 114 members of the Jamaat al Muslimeen, led by Yasin Abu Bakr (formerly known as Lennox Phillip) stormed the Red House (the seat of Parliament), and Trinidad and Tobago Television, the only television station in the country at the time, holding Robinson and country's government hostage for six days before surrendering.[82] The coup leaders were promised amnesty, but upon their surrender they were then arrested, but later released after protracted legal wrangling.[58]
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+
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+ The PNM under Patrick Manning returned to power following the 1991 Trinidad and Tobago general election.[16] Hoping to capitalise on an improvement in the economy, Manning called an early election in 1995, however, this resulted in a hung parliament. Two NAR representatives backed the opposition United National Congress (UNC), which had split off from the NAR in 1989, and they thus took power under Basdeo Panday, who became the country's first Indo-Trinidadian Prime Minister.[16][80][83] After a period of political confusion caused by a series of inconclusive election results, Patrick Manning returned to power in 2001, retaining that position until 2010.[16]
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+
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+ Since 2003 the country entered a second oil boom, and petroleum, petrochemicals and natural gas continue to be the backbone of the economy. Tourism and the public service are the mainstay of the economy of Tobago, though authorities have attempted to diversify the island's economy.[84] A corruption scandal resulted in Manning's defeat by the newly formed People's Partnership coalition in 2010, with Kamla Persad-Bissessar becoming the country's first female Prime Minister.[85][86][87] However, corruption allegations bedevilled the new administration, and the PP were defeated in 2015 by the PNM under Keith Rowley.[88][89]
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+
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+ Trinidad and Tobago is situated between 10° 2' and 11° 12' N latitude and 60° 30' and 61° 56' W longitude, with the Caribbean Sea to the north, the Atlantic Ocean to the east and south, and the Gulf of Paria to the west. It is located in the far south-east of the Caribbean region, with the island of Trinidad being just 11 kilometres (6.8 mi) off the coast of Venezuela in mainland South America across the Columbus Channel.[16] Covering an area of 5,128 km2 (1,980 sq mi),[90] the country consists of two main islands, Trinidad and Tobago, separated by a 20m (30 km) strait, plus a number of much smaller islands, including Chacachacare, Monos, Huevos, Gaspar Grande (or Gasparee), Little Tobago, and Saint Giles Island.[16]
68
+
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+ Trinidad is 4,768 km2 (1,841 sq mi) in area (comprising 93.0% of the country's total area) with an average length of 80 kilometres (50 mi) and an average width of 59 kilometres (37 mi). Tobago has an area of about 300 km2 (120 sq mi), or 5.8% of the country's area, is 41 km (25 mi) long and 12 km (7.5 mi) at its greatest width. Trinidad and Tobago lie on the continental shelf of South America, and are thus geologically considered to lie entirely in South America.[16]
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+
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+ The terrain of the islands is a mixture of mountains and plains.[15] On Trinidad the Northern Range runs parallel with the north coast, and contains the country's highest peak (El Cerro del Aripo), which is 940 metres (3,080 ft) above sea level[15] and second highest (El Tucuche, 936 metres (3,071 ft)).[16] The rest of the island is generally flatter, excluding the Central Range and Montserrat Hills in the centre of the island and the Southern Range and Trinity Hills in the south. The east coast is noted for its beaches, most notably Manzanilla Beach. The island contains several large swamp areas, such as the Caroni Swamp and the Nariva Swamp.[16] Major bodies of water on Trinidad include the Hollis Reservoir, Navet Reservoir, Caroni Reservoir. Trinidad is made up of a variety of soil types, the majority being fine sands and heavy clays. The alluvial valleys of the Northern Range and the soils of the East–West Corridor are the most fertile.[91][citation needed] Trinidad is also notable for containing Pitch Lake, the largest natural reservoir of asphalt in the world.[15][16] Tobago contains a flat plain in its south-west, with the eastern half of the island being more mountainous, culminating in Pigeon Peak, the island's highest point at 550 metres (1,800 ft).[92] Tobago also contains several coral reefs off its coast.[16]
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+
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+ The majority of the population reside on the island of Trinidad, and this is thus the location of largest towns and cities. There are four major municipalities in Trinidad: the capital Port of Spain, San Fernando, Arima and Chaguanas. The main town on Tobago is Scarborough.
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+
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+ The Northern Range consists mainly of Upper Jurassic and Cretaceous metamorphic rocks. The Northern Lowlands (the East–West Corridor and Caroni Plain) consist of younger shallow marine clastic sediments. South of this, the Central Range fold and thrust belt consists of Cretaceous and Eocene sedimentary rocks, with Miocene formations along the southern and eastern flanks. The Naparima Plain and the Nariva Swamp form the southern shoulder of this uplift.[citation needed]
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+
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+ The Southern Lowlands consist of Miocene and Pliocene sands, clays, and gravels. These overlie oil and natural gas deposits, especially north of the Los Bajos Fault. The Southern Range forms the third anticlinal uplift. The rocks consist of sandstones, shales, siltstones and clays formed in the Miocene and uplifted in the Pleistocene. Oil sands and mud volcanoes are especially common in this area.[citation needed]
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+
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+ Trinidad and Tobago has a maritime tropical climate.[15][16] There are two seasons annually: the dry season for the first five months of the year, and the rainy season in the remaining seven of the year. Winds are predominantly from the northeast and are dominated by the northeast trade winds. Unlike many Caribbean islands Trinidad and Tobago lies outside the main hurricane alleys; nevertheless, the island of Tobago was struck by Hurricane Flora on September 30, 1963. In the Northern Range of Trinidad, the climate is often cooler than that of the sweltering heat of the plains below, due to constant cloud and mist cover, and heavy rains in the mountains.
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+
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+ Record temperatures for Trinidad and Tobago are 39 °C (102 °F)[93] for the high in Port of Spain, and a low of 12 °C (54 °F).[94]
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+
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+ Because Trinidad and Tobago lies on the continental shelf of South America, and in ancient times were physically connected to the South American mainland, its biological diversity is unlike that of most other Caribbean islands, and has much more in common with that of Venezuela.[95] The main ecosystems are: coastal and marine (coral reefs, mangrove swamps, open ocean and seagrass beds); forest; freshwater (rivers and streams); karst; man-made ecosystems (agricultural land, freshwater dams, secondary forest); and savanna. On 1 August 1996, Trinidad and Tobago ratified the 1992 Rio Convention on Biological Diversity, and it has produced a biodiversity action plan and four reports describing the country's contribution to biodiversity conservation. These reports formally acknowledged the importance of biodiversity to the well-being of the country's people through provision of ecosystem services.[96]
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+
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+ Information about vertebrates is good, with 472 bird species (2 endemics), about 100 mammals, about 90 reptiles (a few endemics), about 30 amphibians (including several endemics), 50 freshwater fish and at least 950 marine fish.[97] Notable mammal species include the ocelot, manatee, collared peccary (known as the quenk locally), agouti, lappe, red brocket deer, otter, weeper capuchin and red howler monkey; there are also some 70 species of bat, including the vampire bat and fringe-lipped bat.[16][98] Amongst the reptiles, the spectacled caiman is the largest, sometimes growing up to 3m.[95] There are also 47 species of snake, including four venomous species, lizards such as the gecko, iguana, matte lizard and also several species of turtle.[16][99] are present. Of the amphibians, the golden tree frog is endemic to Trinidad.[99] Marine life is abundant, with several species of sea urchin, coral, lobster, anemone, starfish, manta ray, dolphin, porpoise and whale shark present in the islands' waters.[100] The introduced lionfish is viewed as a pest, as it eats many native species of fish and has no natural predators; efforts are currently underway to cull the numbers of this species.[100]
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+
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+ Trinidad and Tobago is noted particularly for its large number of bird species, and is a popular destination for bird watchers. Notable species include the scarlet ibis, cocrico, egret, shiny cowbird, bananaquit, oilbird and various species of honeycreeper, trogon, toucan, parrot, tanager, woodpecker, antbird, kites, hawks, boobies, pelicans and vultures; there are also 17 species of hummingbird, including the tufted coquette which is the world's third smallest.[101]
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+
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+ Information about invertebrates is dispersed and very incomplete. About 650 butterflies,[97] at least 672 beetles (from Tobago alone)[102] and 40 corals[97] have been recorded.[97] Other notable invertebrates include the cockroach, leaf-cutter ant and numerous species of mosquitoes, termites, spiders and tarantulas.
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+
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+ Although the list is far from complete, 1,647 species of fungi, including lichens, have been recorded.[103][104][105] The true total number of fungi is likely to be far higher, given the generally accepted estimate that only about 7% of all fungi worldwide have so far been discovered.[106] A first effort to estimate the number of endemic fungi tentatively listed 407 species.[107]
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+
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+ Information about micro-organisms is dispersed and very incomplete. Nearly 200 species of marine algae have been recorded.[97] The true total number of micro-organism species must be much higher.
94
+
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+ Thanks to a recently published checklist, plant diversity in Trinidad and Tobago is well documented with about 3,300 species (59 endemic) recorded.[97] Despite significant felling, forests still cover about 40% of the country, and there are about 350 different species of tree.[95] A notable tree is the manchineel which is extremely poisonous to humans, and even just touching its sap can cause severe blistering of the skin; the tree is often covered with warning signs.
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+
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+ Trinidad and Tobago is a republic with a two-party system and a bicameral parliamentary system based on the Westminster System.[15]
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+
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+ The head of state of Trinidad and Tobago is the President, currently Paula Mae Weekes.[15] This largely ceremonial role replaced that of the Governor-General (representing the Monarch of Trinidad and Tobago) upon Trinidad and Tobago's becoming a republic in 1976.[16] The head of government is the Prime Minister, currently Keith Rowley.[15] The President is elected by an Electoral college consisting of the full membership of both houses of Parliament.
100
+ The Prime Minister is elected following a general election which takes place every five years. The President is required to appoint the leader of the party who in his or her opinion has the most support of the members of the House of Representatives to this post; this has generally been the leader of the party which won the most seats in the previous election (except in the case of the 2001 General Elections).[16]
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+
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+ Since 1980 Tobago has also had its own elections, separate from the general elections. In these elections, members are elected and serve in the unicameral Tobago House of Assembly.[108][15][16]
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+
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+ Parliament consists of the Senate (31 seats) and the House of Representatives (41 seats, plus the Speaker).[109][15] The members of the Senate are appointed by the president; 16 Government Senators are appointed on the advice of the Prime Minister, six Opposition Senators are appointed on the advice of the Leader of the Opposition, currently Kamla Persad-Bissessar, and nine Independent Senators are appointed by the President to represent other sectors of civil society. The 41 members of the House of Representatives are elected by the people for a maximum term of five years in a "first past the post" system.
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+
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+ Trinidad is split into 14 regional corporations and municipalities, consisting of nine regions and five municipalities, which have a limited level of autonomy.[15][16] The various councils are made up of a mixture of elected and appointed members. Elections are held every three years.[citation needed] The country was formerly divided into counties.
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+
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+ The two main parties are the People's National Movement (PNM) and the United National Congress (UNC); another recent party was the Congress of the People (COP). Support for these parties appears to fall along ethnic lines, with the PNM consistently obtaining a majority of Afro-Trinidadian vote, and the UNC gaining a majority of Indo-Trinidadian support.
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+
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+ The Trinidad and Tobago Defence Force (TTDF) is the military organisation responsible for the defence of the twin island Republic of Trinidad and Tobago.[15] It consists of the Regiment, the Coast Guard, the Air Guard and the Defence Force Reserves. Established in 1962 after Trinidad and Tobago's independence from the United Kingdom, the TTDF is one of the largest military forces in the Anglophone Caribbean.[citation needed]
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+
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+ Its mission statement is to "defend the sovereign good of The Republic of Trinidad and Tobago, contribute to the development of the national community and support the State in the fulfillment of its national and international objectives". The Defence Force has been engaged in domestic incidents, such as the 1990 Coup Attempt, and international missions, such as the United Nations Mission in Haiti between 1993 and 1996.
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+
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+ In 2019, Trinidad and Tobago signed the UN treaty on the Prohibition of Nuclear Weapons.[110]
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+
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+ Trinidad and Tobago maintains close relations with its Caribbean neighbours and major North American and European trading partners. As the most industrialised and second-largest country in the Anglophone Caribbean, Trinidad and Tobago has taken a leading role in the Caribbean Community (CARICOM), and strongly supports CARICOM economic integration efforts. It also is active in the Summit of the Americas process and supports the establishment of the Free Trade Area of the Americas, lobbying other nations for seating the Secretariat in Port of Spain.[citation needed]
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+
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+ As a member of CARICOM, Trinidad and Tobago strongly backed efforts by the United States to bring political stability to Haiti, contributing personnel to the Multinational Force in 1994. After its 1962 independence, Trinidad and Tobago joined the United Nations and Commonwealth of Nations. In 1967 it became the first Commonwealth country to join the Organization of American States (OAS).[111] In 1995 Trinidad played host to the inaugural meeting of the Association of Caribbean States and has become the seat of this 35-member grouping, which seeks to further economic progress and integration among its states. In international forums, Trinidad and Tobago has defined itself as having an independent voting record, but often supports US and EU positions.[citation needed]
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+
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+ Trinidad and Tobago has in recent decades suffered from a relatively high crime rate;[112][113] there are currently roughly 500 murders per year.[114][76] The country is a noted transshipment centre for the trafficking of illegal drugs from South America to the rest of the Caribbean and beyond to North America.[115] Some estimates put the size of the 'hidden economy' as high as 20–30% of measured GDP.[116]
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+
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+ Though there have been no terrorism-related incidents in the country since the 1990 Islamic coup attempt, Trinidad and Tobago remains a potential target; for example, in February 2018 a plan to attack the Carnival was foiled by police.[113] It is estimated that roughly 100 citizens of the country have traveled to the Middle East to fight for Islamic State.[112][113] In 2017 the government adopted a counter-terrorism and extremism strategy.[113]
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+
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+ The country's prison administration is the Trinidad and Tobago Prison Service (TTPS), it is under the control of the Commissioner of Prisons Gerard Wilson, located in Port-of-Spain.[117] The prison population rate is 292 people per 100,000. The total prison population, including pre-trial detainees and remand prisoners, is 3,999 prisoners. The population rate of pre-trial detainees and remand prisoners is 174 per 100,000 of the national population (59.7% of the prison population). In 2018, the female prison population rate is 8.5 per 100,000 of the national population (2.9% of the prison population). Prisoners that are minors makes up 1.9% of the prison population and foreigners prisoners make 0.8% of the prison population. The occupancy level of Trinidad and Tobago's prison system is at 81.8% capacity.[117] Trinidad and Tobago has nine prison establishments; Golden Grove Prison, Maximum Security Prison, Port of Spain Prison, Eastern Correctional Rehabilitation Centre, Remand Prison, Tobago Convict Prison, Carrera Convict Island Prison, Women's Prison and Youth Training and Rehabilitation Centre.[118] Trinidad and Tobago also use labor yards as prisons, or means of punishment.[119]
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+
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+ The population of the country currently stands at 1,363,985 (July 2019 est.).[citation needed]
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+
128
+ The ethnic composition of Trinidad and Tobago reflects a history of conquest and immigration.[120] While the earliest inhabitants were of Amerindian heritage, the two dominant groups in the country are now those of South Asian and of African heritage. Indo-Trinidadian and Tobagonians make up the country's largest ethnic group (approximately 35.4%);[15] they are primarily the descendants of indentured workers from South Asia (mostly from India), brought to replace freed African slaves who refused to continue working on the sugar plantations. Through cultural preservation many residents of Indian descent continue to maintain traditions from their ancestral homeland. Indo-Trinidadians reside primarily on Trinidad; as of the 2011 census only 2.5% of Tobago's population was of Indian descent.[121]
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+
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+ Afro-Trinidadians and Tobagonians make up the country's second largest ethnic group, with approximately 34.2% of the population identifying as being of African descent.[15] The majority of people of an African background are the descendants of slaves forcibly transported to the islands from as early as the 16th century. This group constitute the majority on Tobago, at 85.2%.[121]
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+
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+ The bulk of the rest of the population are those who identify as being of mixed heritage.[15] There are also small but significant minorities of people of Amerindian, European, Chinese, and Arab descent. Arima on Trinidad is a noted centre of Amerindian culture.[16]
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+
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+ English is the country's official language (the local variety of standard English is Trinidadian and Tobagonian English or more properly, Trinidad and Tobago Standard English, abbreviated as "TTSE"), but the main spoken language is either of two English-based creole languages (Trinidadian Creole or Tobagonian Creole), which reflects the Amerindian, European, African, and Asian heritage of the nation. Both creoles contain elements from a variety of African languages; Trinidadian English Creole, however, is also influenced by French and French Creole (Patois).[122]
135
+
136
+ A majority of the early Indian immigrants spoke the Bhojpuri and Awadhi dialect of Hindustani (Hindi-Urdu), which later formed into Trinidadian Hindustani (Hindi-Urdu), which became the lingua franca of Indo-Trinidadian and Tobagonians. From 1935 Indian films began showing to audiences in Trinidad; most of these were in the Standard Hindustani (Hindi-Urdu) dialect and this modified Trinidadian Hindustani slightly by adding Standard Hindu and Urdu phrases and vocabulary to Trinidadian Hindustani. Indian films also revitalised Hindustani among Indo-Trinidadian and Tobagonians.[123] Around the mid to late 1970s the lingua franca of Indo-Trinidaian and Tobagonians switched from Trinidadian Hindustani to a sort of Hindinised version of English. Today Hindustani survives on through Indo-Trinidadian and Tobagonian musical forms such as, Bhajan, Indian classical music, Indian folk music, Filmi, Pichakaree, Chutney, Chutney soca, and Chutney parang. Presently there are about 26,000 people, which is 5.53% of the Indo-Trinidadian and Tobagonian population, who speak Trinidadian Hindustani. Many Indo-Trinidadians and Tobagonians today speak a type of Hinglish that consist of Trinidadian and Tobagonian English that is heavily laced with Trinidadian Hindustani vocabulary and phrases and many Indo-Trinidadians and Tobagonians can recite phrases or prayers in Hindustani today. There are many places in Trinidad and Tobago that have names of Hindustani origin. Some phrases and vocabulary have even made their way into the mainstream English and English Creole dialects of the country.[124][125][126][127][128]
137
+
138
+ The Chinese language first came to Trinidad and Tobago in 1806, when the British had brought Chinese labourers in order to determine if they were fit to use as labourers after the abolition of slavery.[citation needed] About 2,645 Chinese immigrants arrived in Trinidad as indentured labour between 1853 and 1866.[citation needed] A majority of the people who immigrated in the 19th century were from southern China and spoke the Hakka and Yue dialects of Chinese. In the 20th century after the years of indentureship up to the present-day more Chinese people have immigrated to Trinidad and Tobago for business and they speak the dialects of the indenturees along with other Chinese dialects, such as Mandarin and Min.[125][129] J. Dyer Ball, writing in 1906, says: "In Trinidad there were, about twenty years ago, 4,000 or 5,000 Chinese, but they have decreased to probably about 2,000 or 3,000, [2,200 in 1900]. They used to work in sugar plantations, but are now principally shopkeepers, as well as general merchants, miners and railway builders,
139
+ etc."[130]
140
+
141
+ The indigenous languages were Yao on Trinidad and Karina on Tobago, both Cariban, and Shebaya on Trinidad, which was Arawakan.[125]
142
+
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+ According to the 2011 census,[3] Roman Catholics were the largest single religious group in Trinidad and Tobago with 21.60% of the total population. The Pentecostal/Evangelical/Full Gospel denominations were the third largest group with 12.02% of the population. The remaining population is made of various Christian denominations (Spiritual Shouter Baptists (5.67%), Anglicans (5.67%), Seventh-day Adventists (4.09%), Presbyterians or Congregationalists (2.49%), Jehovah's Witnesses (1.47%), other Baptists (1.21%), Methodists (0.65%) and the Moravian Church (0.27%)). Respondents who did not state a religious affiliation represented 11.1% of the population, with 2.18% declaring themselves Irreligious.
144
+
145
+ Hindus were the second largest group with 18.15%.[3] Hinduism is practiced throughout the country and Diwali is a public holiday, and other Hindu holidays are also widely celebrated.
146
+
147
+ Muslims represent 4.97% of the population.[3] Eid al-Fitr is a public holiday and Eid al-Adha, Mawlid, Hosay, and other Muslim holidays are also celebrated. There has also been a Jewish community on the islands for many centuries, however their numbers have never been large, with a 2007 estimating putting the Jewish population at 55 individuals.[131][132]
148
+
149
+ African-derived or Afrocentric religions are also practised, notably Trinidad Orisha (Yoruba) believers (0.9%) and Rastafarians (0.27%).[3] Various aspects of traditional obeah beliefs are still commonly practised on the islands.[50]
150
+
151
+ Two African syncretic faiths, the Shouter or Spiritual Baptists and the Orisha faith (formerly called Shangos, a less than complimentary term)[citation needed] are among the fastest growing religious groups. Similarly, there is a noticeable increase in numbers of Evangelical Protestant and Fundamentalist churches usually lumped as "Pentecostal" by most Trinidadians, although this designation is often inaccurate. Sikhism, Jainism, Bahá'í, and Buddhism are practised by a minority of Indo-Trinidadian and Tobagonians. Several eastern religions such as Buddhism and Chinese folk religions such as Taoism and Confucianism are followed by Chinese Trinidadian and Tobagonian.
152
+
153
+ Children generally start pre-school at two and a half years but this is not mandatory. They are however, expected to have basic reading and writing skills when they commence primary school. Students proceed to a primary school at the age of five years. Seven years are spent in primary school. The seven classes of primary school consists of First Year and Second Year, followed by Standard One through Standard Five. During the final year of primary school, students prepare for and sit the Secondary Entrance Assessment (SEA) which determines the secondary school the child will attend.
154
+
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+ Students attend secondary school for a minimum of five years, leading to the CSEC (Caribbean Secondary Education Certificate) examinations, which is the equivalent of the British GCSE O levels. Children with satisfactory grades may opt to continue high school for a further two-year period, leading to the Caribbean Advanced Proficiency Examinations (CAPE), the equivalent of GCE A levels. Both CSEC and CAPE examinations are held by the Caribbean Examinations Council (CXC). Public Primary and Secondary education is free for all, although private and religious schooling is available for a fee.
156
+
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+ Tertiary education for tuition costs are provided for via GATE (The Government Assistance for Tuition Expenses), up to the level of the bachelor's degree, at the University of the West Indies (UWI), the University of Trinidad and Tobago (UTT), the University of the Southern Caribbean (USC), the College of Science, Technology and Applied Arts of Trinidad and Tobago (COSTAATT) and certain other local accredited institutions. Government also currently subsidises some Masters programmes. Both the Government and the private sector also provide financial assistance in the form of academic scholarships to gifted or needy students for study at local, regional or international universities.
158
+
159
+ While women account for only 49% of the population, they constitute nearly 55% of the workforce in the country.[137]
160
+
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+ Trinidad and Tobago is the most developed nation and one of the wealthiest in the Caribbean and is listed in the top 40 (2010 information) of the 70 high-income countries in the world.[citation needed] Its gross national income per capita of US$20,070[138] (2014 gross national income at Atlas Method) is one of the highest in the Caribbean.[139] In November 2011, the OECD removed Trinidad and Tobago from its list of developing countries.[140] Trinidad's economy is strongly influenced by the petroleum industry. Tourism and manufacturing are also important to the local economy. Tourism is a growing sector, particular on Tobago, although proportionately it is much less important than in many other Caribbean islands. Agricultural products include citrus and cocoa. It also supplies manufactured goods, notably food, beverages, and cement, to the Caribbean region.
162
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+ Trinidad and Tobago is the leading Caribbean producer of oil and gas, and its economy is heavily dependent upon these resources.[16] Oil and gas account for about 40% of GDP and 80% of exports, but only 5% of employment.[15] Recent growth has been fuelled by investments in liquefied natural gas (LNG), petrochemicals, and steel. Additional petrochemical, aluminium, and plastics projects are in various stages of planning.
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+ The country is also a regional financial centre, and the economy has a growing trade surplus.[90] The expansion of Atlantic LNG over the past six years created the largest single-sustained phase of economic growth in Trinidad and Tobago. The nation is an exporter of LNG and supplied a total of 13.4 billion m3 in 2017. The largest markets for Trinidad and Tobago's LNG exports are Chile and the United States.[141]
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+ Trinidad and Tobago has transitioned from an oil-based economy to a natural gas based economy. In 2017, natural gas production totalled 18.5 billion m3, a decrease of 0.4% from 2016 with 18.6 billion m3 of production.[141] Oil production has decreased over the past decade from 7.1 million metric tonnes per year in 2007 to 4.4 million metric tonnes per year in 2017.[142] In December 2005, the Atlantic LNG's fourth production module or "train" for liquefied natural gas (LNG) began production. Train four has increased Atlantic LNG's overall output capacity by almost 50% and is the largest LNG train in the world at 5.2 million tons/year of LNG.[citation needed]
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+ Trinidad and Tobago is far less dependent on tourism than many other Caribbean countries and territories, with the bulk of tourist activity occurring on Tobago.[16] The government has made efforts to boost this sector in recent years.[16]
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+ Historically agricultural production (for example, sugar and coffee) dominated the economy, however this sector has been in steep decline since the 20th century and now forms just 0.4% of the country's GDP, employing 3.1% of the workforce.[15][16] Various fruits and vegetables are grown, such as cucumbers, eggplant, cassava, pumpkin, dasheen (taro) and coconut; fishing is still also commonly practised.[15]
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+ Trinidad and Tobago, in an effort to undergo economic transformation through diversification,[15] formed InvesTT in 2012 to serve as the country's sole investment promotion agency. This agency is aligned to the Ministry of Trade and Industry and is to be the key agent in growing the country's non-oil and gas sectors significantly and sustainably.[143]
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+ Trinidad and Tobago has a well developed communications sector. The telecommunications and broadcasting sectors generated an estimated TT$5.63 billion (US$0.88 billion) in 2014, which as a percentage of GDP equates to 3.1 percent. This represented a 1.9 percent increase in total revenues generated by this industry compared to last year. Of total telecommunications and broadcasting revenues, mobile voice services accounted for the majority of revenues with TT$2.20 billion (39.2 percent). This was followed by internet services which contributed TT$1.18 billion or 21.1 percent. The next highest revenue earners for the industry were fixed voice services and paid television services whose contributions totalled TT$0.76 billion and TT$0.70 billion respectively (13.4 percent and 12.4 percent). International voice services was next in line, generating TT$0.27 billion (4.7 percent) in revenues. Free-to Air radio and television services contributed TT$0.18 billion and TT$0.13 billion respectively (3.2 percent and 2.4 percent). Finally, other contributors included "other revenues" and "leased line services" with earnings of TT$0.16 billion and TT$0.05 billion respectively, with 2.8 percent and 0.9 percent.[144]
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+ There are several providers for each segment of the telecommunications market. Fixed Lines Telephone service is provided by Digicel, TSTT (operating as bmobile) and Cable & Wireless Communications operating as FLOW; cellular service is provided by TSTT (operating as bmobile) and Digicel whilst internet service is provided by TSTT, FLOW, Digicel, Green Dot and Lisa Communications.
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+ The Government of Trinidad and Tobago has recognised the creative industries as a pathway to economic growth and development. It is one of the newest, most dynamic sectors where creativity, knowledge and intangibles serve as the basic productive resource. In 2015, the Trinidad and Tobago Creative Industries Company Limited (CreativeTT) was established as a state agency under the Ministry of Trade and Industry with a mandate to stimulate and facilitate the business development and export activities of the Creative Industries in Trinidad and Tobago to generate national wealth, and, as such, the company is responsible for the strategic and business development of the three (3) niche areas and sub sectors currently under its purview – Music, Film and Fashion. MusicTT, FilmTT and FashionTT are the subsidiaries established to fulfil this mandate.
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+ The transport system in Trinidad and Tobago consists of a dense network of highways and roads across both major islands, ferries connecting Port of Spain with Scarborough and San Fernando, and international airports on both islands.[16] The Uriah Butler Highway, Churchill Roosevelt Highway and the Sir Solomon Hochoy Highway links the island of Trinidad together, whereas the Claude Noel Highway is the only major highway in Tobago. Public transportation options on land are public buses, private taxis and minibuses. By sea, the options are inter-island ferries and inter-city water taxis.[145]
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+ The island of Trinidad is served by Piarco International Airport located in Piarco, which opened on 8 January 1931.[citation needed] Elevated at 17.4 metres (57 ft) above sea level it comprises an area of 680 hectares (1,700 acres) and has a runway of 3,200 metres (10,500 ft). The airport consists of two terminals, the North Terminal and the South Terminal. The older South Terminal underwent renovations in 2009 for use as a VIP entrance point during the 5th Summit of the Americas. The North Terminal was completed in 2001, and consists of[146] 14-second-level aircraft gates with jetways for international flights, two ground-level domestic gates and 82 ticket counter positions.
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+ In 2008 the passenger throughput at Piarco International Airport was approximately 2.6 million. It is the seventh busiest airport in the Caribbean and the third busiest in the English-speaking Caribbean, after Sangster International Airport and Lynden Pindling International Airport.[citation needed] Caribbean Airlines, the national airline, operates its main hub at the Piarco International Airport and services the Caribbean, the United States, Canada and South America. The airline is wholly owned by the Government of Trinidad and Tobago. After an additional cash injection of US$50 million, the Trinidad and Tobago government acquired the Jamaican airline Air Jamaica on 1 May 2010, with a 6–12-month transition period to follow.[147]
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+ The Island of Tobago is served by the A.N.R. Robinson International Airport in Crown Point.[16] This airport has regular services to North America and Europe. There are regular flights between the two islands, with fares being heavily subsidised by the Government.
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+ Trinidad was formerly home to a railway network, however this was closed down in 1968.[148] There have been talks to build a new railway on the islands, though nothing yet has come of this.[149]
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+ The Strategic Plan for the Caribbean Community 2015–2019 was adopted by Trinidad and Tobago and the other members of the Caribbean Common Market (CARICOM) in 2014. The first of its kind, this document reflects a desire among countries to embrace a more profound regionalism, in order to reposition the Caribbean in an increasingly volatile global economy. The plan proposes mobilising funding from the public and private sectors to foster research and development (R&D) in the business sector. The plan outlines strategies for nurturing creativity, entrepreneurship, digital literacy and for making optimum use of available resources. It focuses on developing creative, manufacturing and service industries, with a special focus on tourism initially, natural resources and value-added products, agriculture and fisheries, to reduce dependence on food imports and foster sustainable fisheries, and energy efficiency.[150][151]
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+ Trinidad and Tobago is the region's leading exporter of oil and gas but imports of fossil fuels provided over 90% of the energy consumed by its CARICOM neighbours in 2008. This vulnerability led CARICOM to develop an Energy Policy which was approved in 2013. This policy is accompanied by the CARICOM Sustainable Energy Roadmap and Strategy (C-SERMS). Under the policy, renewable energy sources are to contribute 20% of the total electricity generation mix in member states by 2017, 28% by 2022 and 47% by 2027.[150]
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+ In 2014 Trinidad and Tobago was the third country in the world which emitted the most CO2 per capita after Qatar and Curacao according to the World Bank.[152] On average, each inhabitant produced 34.2 metric tons of CO2 in the atmosphere. In comparison, the world average was 5.0 tons per capita the same year.
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+ The Caribbean Community Climate Change Centre (CCCCC) has produced an implementation plan for 2011–2021 and carried out work to assess and build capacity in climate change mitigation and resilient development strategies. This work has been supported by the region's specialists, who have produced models for climate change and mitigation processes in Caribbean states. They also play a major advisor role to the divisions in ministries responsible for climate change. The growing frequency and intensity of hurricanes is of concern to all Caribbean nations. In 2012, Trinidad and Tobago had a 9% chance each year of being struck by a hurricane, according to estimates by the International Monetary Fund.[150][153]
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+ The two main bodies responsible for science, technology and innovation in Trinidad and Tobago are the Ministry of Science, Technology and Higher Education and the National Commission for Science and Technology.[150]
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+ In 2002, Trinidad and Tobago adopted Vision 2020. Like Jamaica's Vision 2030 (2009) and the Strategic Plan of Barbados for 2005–2025, Trinidad and Tobago's Vision 2020 accords central importance to harnessing science, technology and innovation (STI) to raise living standards and strengthen resilience to environmental shocks like hurricanes.[150]
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+ Trinidad and Tobago is one of the more affluent members of CARICOM, thanks to its natural resources. Despite this, it spent just 0.05% of GDP on R&D in 2012, according to the UNESCO Institute for Statistics. Even when the country was enjoying economic growth of 8% per annum in 2004, it devoted just 0.11% of GDP to R&D. Calculated in thousands of current Purchasing Power Parity (PPP) dollars, research expenditure actually dropped between 2009 and 2012 from 21 309 to 19 232. This corresponds to research expenditure of $PPP 65 per capita in 2009 and $PPP 45 in 2012.[150]
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+ Industrial R&D has declined since 2000, perhaps owing to the drop in research activity in the sugar sector. Whereas industrial R&D accounted for 24% of domestic research expenditure in 2004 and 29.5% in 2005, it had become almost non-existent by 2010.[150]
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+ The number of researchers in Trinidad and Tobago grew from 787 to 914 between 2009 and 2012. This corresponds in a rise from 595 to 683 in the number of researchers (head counts) per million inhabitants.[150]
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+ Scientific output grew between 2007 and 2011, according to Thomson Reuters' Web of Science (Science Citation Index Expanded) before contracting over the period 2012–2014. Trinidad and Tobago produced 109 publications per million population in 2014, behind Grenada (1,430), St Kitts and Nevis (730), Barbados (182) and Dominica (138) but ahead of the Bahamas (86), Belize (47) and Jamaica (42).[150]
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+ Between 2008 and 2014, scientists collaborated most with their peers from the United States (251 papers), United Kingdom (183), Canada (95), India (63) and Jamaica (43), according to the copublication record of Thomson Reuters. In turn, Jamaican scientists considered their counterparts from Trinidad and Tobago to be their fourth-closest collaborators (with 43 joint papers) after those from the United States, United Kingdom and Canada.[150]
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+ Between 2008 and 2013, Trinidad and Tobago registered 17 patents with the US Patent and Trademark Office (USPTO). This corresponds to 13% of the 134 patents registered by CARICOM members over this period. The top contributors were the Bahamas (34 patents) and Jamaica (22).[150]
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+ Trinidad and Tobago led CARICOM members for the value of high-tech exports in 2008 (US$36.2 million) but these exports plummeted to US$3.5 million the following year, according to the Comtrade database of the United Nations Statistics Division.[150]
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+ The Caribbean Industrial Research Institute in Trinidad and Tobago facilitates climate change research and provides industrial support for R&D related to food security. It also carries out equipment testing and calibration for major industries.[150]
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+ The Caribbean Epidemiology Centre in Port of Spain, University of Trinidad and Tobago, Tobago Institute of Health, and University of the West Indies (St Augustine campus) also conduct R&D.[150]
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+ Probability of a hurricane striking Caribbean countries in a given year, 2012 (%).[154]
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+ Electricity costs for the CARICOM countries, 2011[155]
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+ GERD by sector of performance in Trinidad and Tobago, 2000–2012[156]
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+ Scientific publication trends in the CARICOM countries, 2005–2014[157]
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+ Scientific publications in the CARICOM countries, 2014[157]
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+ USPTO patents granted to Caribbean countries, 2008–2013.[158]
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+ Trinidad and Tobago has a diverse culture mixing Indian, African, Creole, Chinese, Amerindian, Arab, Latino, and European influences, reflecting the various communities who have migrated to the islands over the centuries. The island is particularly renowned for its annual Carnival celebrations.[16] Festivals rooted in various religions and cultures practiced on the islands are also popular, such as Christmas, Divali, Phagwah (Holi), Easter, New Year’s Day, Hosay, Eid al-Fitr, the Santa Rosa Indigenous Festival, and Chinese New Year.
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+ Trinidad and Tobago claims two Nobel Prize-winning authors, V. S. Naipaul and St Lucian-born Derek Walcott (who also founded the Trinidad Theatre Workshop). Other notable writers include Neil Bissoondath, Vahni Capildeo, Earl Lovelace, Seepersad Naipaul, Shiva Naipaul, Lakshmi Persaud, Kenneth Ramchand, Arnold Rampersad, and Samuel Selvon.
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+ Trinidadian designer Peter Minshall is renowned not only for his Carnival costumes but also for his role in opening ceremonies of the Barcelona Olympics, the 1994 FIFA World Cup, the 1996 Summer Olympics, and the 2002 Winter Olympics, for which he won an Emmy Award.[161]
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+ Trinidad and Tobago is the birthplace of calypso music and the steelpan.[162][163][164] Trinidad is also the birthplace of soca music, chutney music, chutney-soca, parang, rapso, pichakaree and chutney parang.
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+ The limbo dance originated in Trinidad as an event that took place at wakes in Trinidad. The limbo has African roots. It was popularized in the 1950s by dance pioneer Julia Edwards[165] (known as the First Lady of Limbo) and her company which appeared in several films.[166] Bélé, Bongo, and whining are also dance forms with African roots.[167]
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+ Jazz, ballroom, ballet, modern, and salsa dancing are also popular.[167]
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+ Indian dance forms are also popular in Trinidad and Tobago.[168] Kathak, Odissi, and Bharatanatyam are the most popular Indian classical dance forms in Trinidad and Tobago.[169] Indian folk dances and Bollywood dances are also popular.[169]
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+ Geoffrey Holder (brother of Boscoe Holder) and Heather Headley are two Trinidad-born artists who have won Tony Awards for theatre. Holder also has a distinguished film career, and Headley has won a Grammy Award as well.
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+ Indian theatre is also popular throughout Trinidad and Tobago. Dramas such as Nautankis, Raja Harishchandra, Sharwan Kumar, and Alha-Khand were brought by Indians to Trinidad and Tobago, however they had largely began to die out, till preservation began by Indian cultural groups.[170] The drama about the life of the Hindu god Rama, Ramleela, is popular during the time between Sharad Navaratri and Dushera and the drama about the life of the Hindu god Krishna, Ras leela (Krishna leela), is popular around the time of Krishna Janmashtami.[171][172][173]
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+ Trinidad and Tobago is also smallest country to have two Miss Universe titleholders and the first black woman ever to win: Janelle Commissiong in 1977, followed by Wendy Fitzwilliam in 1998; the country has also had one Miss World titleholder, Giselle LaRonde.[citation needed]
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+ Hasely Crawford won the first Olympic gold medal for Trinidad and Tobago in the men's 100-metre dash in the 1976 Summer Olympics. Nine different athletes from Trinidad and Tobago have won twelve medals at the Olympics, beginning with a silver medal in weightlifting, won by Rodney Wilkes in 1948,[174] and most recently, a gold medal by Keshorn Walcott in the men's javelin throw in 2012. Ato Boldon has won the most Olympic and World Championship medals for Trinidad and Tobago in athletics, with eight in total – four from the Olympics and four from the World Championships. Boldon was the sole world champion Trinidad and Tobago has produced until Jehue Gordon in Moscow 2013. Ato won the 1997 200 m sprint World Championship in Athens. Swimmer George Bovell III won a bronze medal in the men's 200 m IM in 2004. At the 2017 World Championship in London, the Men 4x400 relay team captured the title, thus the country now celebrates three world championships titles. The team consisted of Jarrin Solomon, Jareem Richards, Machel Cedenio and Lalonde Gordon with Renny Quow who ran in the heats.
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+ Also in 2012 Lalonde Gordon competed in the XXX Summer Olympics where he won a bronze medal in the 400 metres (1,300 feet), being surpassed by Luguelin Santos of the Dominican Republic and Kirani James of Grenada. Keshorn Walcott (as stated above) came first in javelin and earned a gold medal, making him the second Trinidadian in the country's history to receive one. This also makes him the first Western[clarification needed] athlete in 40 years to receive a gold medal in the javelin sport, and the first athlete from Trinidad and Tobago to win a gold medal in a field event in the Olympics. Sprinter Richard Thompson is also from Trinidad and Tobago. He came second place to Usain Bolt in the Beijing Olympics in the 100 metres (330 feet) with a time of 9.89s.
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+ In 2018 The Court of Arbitration for Sport made its final decision on the failed doping sample from the Jamaican team in the 4 x 100 relay in the 2008 Olympic Games. The team from Trinidad and Tobago will be awarded the gold medal, because of the second rank during the relay run.[175]
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+ Cricket is a popular sport of Trinidad and Tobago, often deemed the national sport, and there is intense inter-island rivalry with its Caribbean neighbours. Trinidad and Tobago is represented at Test cricket, One Day International as well as Twenty20 cricket level as a member of the West Indies team. The national team plays at the first-class level in regional competitions such as the Regional Four Day Competition and Regional Super50. Meanwhile, the Trinbago Knight Riders play in the Caribbean Premier League.
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+ The Queen's Park Oval located in Port of Spain is the largest cricket ground in the West Indies, having hosted 60 Test matches as of January 2018. Trinidad and Tobago along with other islands from the Caribbean co-hosted the 2007 Cricket World Cup.
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+ Brian Lara, world record holder for the most runs scored both in a Test and in a First Class innings and other records, was born in a small town of Santa Cruz and is often referred to as the Prince of Port of Spain or simply the Prince. This legendary West Indian batsman is widely regarded (along with Sir Donald Bradman, Sunil Gavaskar and Sachin Tendulkar[citation needed]) as one of the best batsmen ever to have played the game,[citation needed] and is one of the most famous sporting icons in the country.
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+ Association football is also a popular sport in Trinidad and Tobago. The men's national football team qualified for the 2006 FIFA World Cup for the first time by beating Bahrain in Manama on 16 November 2005, making them the second smallest country ever (in terms of population) to qualify, after Iceland. The team, coached by Dutchman Leo Beenhakker, and led by Tobagonian-born captain Dwight Yorke, drew their first group game – against Sweden in Dortmund, 0–0, but lost the second game to England on late goals, 0–2. They were eliminated after losing 2–0 to Paraguay in the last game of the Group stage. Prior to the 2006 World Cup qualification, Trinidad and Tobago came close in a controversial qualification campaign for the 1974 FIFA World Cup. Following the match, the referee of their critical game against Haiti was awarded a lifetime ban for his actions.[176] Trinidad and Tobago again fell just short of qualifying for the World Cup in 1990, needing only a draw at home against the United States but losing 1–0.[177] They play their home matches at the Hasely Crawford Stadium. Trinidad and Tobago hosted the 2001 FIFA U-17 World Championship, and hosted the 2010 FIFA U-17 Women's World Cup.
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+ The TT Pro League is the country's primary football competition and is the top level of the Trinidad and Tobago football league system. The Pro League serves as a league for professional football clubs in Trinidad and Tobago. The league began in 1999 as part of a need for a professional league to strengthen the country's national team and improve the development of domestic players. The first season took place in the same year beginning with eight teams.
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+ Basketball is commonly played in Trinidad and Tobago in colleges, universities and throughout various urban basketball courts. Its national team is one of the most successful teams in the Caribbean. At the Caribbean Basketball Championship it won four straight gold medals from 1986 to 1990.
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+ Netball has long been a popular sport in Trinidad and Tobago, although it has declined in popularity in recent years. At the Netball World Championships they co-won the event in 1979, were runners up in 1987, and second runners up in 1983.
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+ Rugby is played in Trinidad and Tobago and continues to be a popular sport, and horse racing is regularly followed in the country.
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+ There is also the Trinidad and Tobago national baseball team which is controlled by the Baseball/Softball Association of Trinidad and Tobago, and represents the nation in international competitions. The team is a provisional member of the Pan American Baseball Confederation.
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+ There are a number of 9 and 18-hole golf courses on Trinidad and Tobago. The most established is the St Andrews Golf Club, Maraval in Trinidad (commonly referred to as Moka), and there is a newer course at Trincity, near Piarco Airport called Millennium Lakes. There are 18-hole courses at Chaguramas and Point-a-Pierre and 9-hole courses at Couva and St Madeline. Tobago has two 18-hole courses. The older of the two is at Mount Irvine, with the Magdalena Hotel & Golf Club (formerly Tobago Plantations) being built more recently.
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+ Although a minor sport, bodybuilding is of growing interest in Trinidad and Tobago. Heavyweight female bodybuilder Kashma Maharaj is of Trinidadian descent.
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+ Dragonboat is also another water-sport that has been rapidly growing over the years. Introduced in 2006. the fraternity made consistent strides in having more members apart of the TTDBF (Trindad and Tobago Dragonboat Federation) as well as performing on an international level such as the 10th IDBF World Nations Dragon Boat Championships in Tampa, Florida in the US in 2011.
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+ Claude Noel is a former world champion in professional boxing. He was born in Tobago.
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+ The flag was chosen by the Independence committee in 1962. Red, black and white symbolise the warmth of the people, the richness of the earth and water respectively.[178][179]
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+ The coat of arms was designed by the Independence committee, and features the scarlet ibis (native to Trinidad), the cocrico (native to Tobago) and hummingbird. The shield bears three ships, representing both the Trinity, and the three ships that Columbus sailed.[178]
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+ There are five categories and thirteen classes of national awards:[180]
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+ The national anthem of the twin-island state is "Forged from the Love of Liberty".[181][182]
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+ Other national songs include "God Bless Our Nation"[183] and "Our Nation's Dawning".[184]
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+ The national flower of Trinidad and Tobago is the chaconia flower. It was chosen as the national flower because it is an indigenous flower that has witnessed the history of Trinidad and Tobago. It was also chosen as the national flower because of its red colour that resembles the red of the national flag and coat of arms and because it blooms around the Independence Day of Trinidad and Tobago.[185]
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+ The national birds of Trinidad and Tobago are the scarlet ibis and the cocrico. The scarlet ibis is kept safe by the government by living in the Caroni Bird Sanctuary which was set up by the government for the protection of these birds. The Cocrico is more indigenous to the island of Tobago and are more likely to be seen in the forest.[186] The hummingbird is considered another symbol of Trinidad and Tobago due to its significance to the indigenous peoples, however, it is not a national bird.[187]
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+ The scarlet ibis birds flying over the Caroni Swamp
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+ The cocrico bird in Tobago
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+ This article incorporates text from a free content work. UNESCO Science Report: towards 2030, 156–173, Harold Ramkissoon & Ishenkumba A. Kahwa, UNESCO Publishing. To learn how to add open license text to Wikipedia articles, please see this how-to page. For information on reusing text from Wikipedia, please see the terms of use.
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+ ^ These three Dutch Caribbean territories form the SSS islands.
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+ * These three Dutch Caribbean territories form the BES islands.
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+ † Physiographically, these are continental islands not a part of the volcanic Windward Islands arc. However, based on proximity, these islands are sometimes grouped with the Windward Islands culturally and politically.
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+ ~ Disputed territories administered by Colombia.
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+ # Physiographically, Bermuda is an isolated oceanic island in the North Atlantic Ocean, not a part of the Antilles, West Indies, Caribbean, North American continent or South American continent. Usually grouped with Northern American countries based on proximity; occasionally grouped with the Caribbean region culturally.
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+ Coordinates: 10°36′N 61°6′W / 10.600°N 61.100°W / 10.600; -61.100
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1
+ In trade, barter (derived from baretor[1]) is a system of exchange where participants in a transaction directly exchange goods or services for other goods or services without using a medium of exchange, such as money.[2] Economists distinguish barter from gift economies in many ways; barter, for example, features immediate reciprocal exchange, not delayed in time. Barter usually takes place on a bilateral basis, but may be multilateral (i.e., mediated through a trade exchange). In most developed countries, barter usually only exists parallel to monetary systems to a very limited extent. Market actors use barter as a replacement for money as the method of exchange in times of monetary crisis, such as when currency becomes unstable (e.g., hyperinflation or a deflationary spiral) or simply unavailable for conducting commerce.
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+ No ethnographic studies have shown that any present or past society has used barter without any other medium of exchange or measurement, nor have anthropologists found evidence that money emerged from barter, instead finding that gift-giving (credit extended on a personal basis with an inter-personal balance maintained over the long term) was the most usual means of exchange of goods and services. Nevertheless, economists since the times of Adam Smith (1723–1790), taking non-specific, often wholely or inaccurately imagined pre-modern societies as examples, have used the inefficiency of barter to explain the emergence of money, of "the" economy, and hence of the discipline of economics itself.[3]
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+ [4][5]
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+ Adam Smith, the father of modern economics, sought to demonstrate that markets (and economies) pre-existed the state. He argued (against conventional wisdom) that money was not the creation of governments. Markets emerged, in his view, out of the division of labor, by which individuals began to specialize in specific crafts and hence had to depend on others for subsistence goods. These goods were first exchanged by barter. Specialization depended on trade, but was hindered by the "double coincidence of wants" which barter requires, i.e., for the exchange to occur, each participant must want what the other has. To complete this hypothetical history, craftsmen would stockpile one particular good, be it salt or metal, that they thought no one would refuse. This is the origin of money according to Smith. Money, as a universally desired medium of exchange, allows each half of the transaction to be separated.[3]
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+ Barter is characterized in Adam Smith's "The Wealth of Nations" by a disparaging vocabulary: "haggling, swapping, dickering." It has also been characterized as negative reciprocity, or "selfish profiteering."[6]
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+ Anthropologists have argued, in contrast, "that when something resembling barter does occur in stateless societies it is almost always between strangers."[7] Barter occurred between strangers, not fellow villagers, and hence cannot be used to naturalistically explain the origin of money without the state. Since most people engaged in trade knew each other, exchange was fostered through the extension of credit.[8][9] Marcel Mauss, author of 'The Gift', argued that the first economic contracts were to not act in one's economic self-interest, and that before money, exchange was fostered through the processes of reciprocity and redistribution, not barter.[10] Everyday exchange relations in such societies are characterized by generalized reciprocity, or a non-calculative familial "communism" where each takes according to their needs, and gives as they have.[11]
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+ Since direct barter does not require payment in money, it can be utilized when money is in short supply, when there is little information about the credit worthiness of trade partners, or when there is a lack of trust between those trading.
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+ Barter is an option to those who cannot afford to store their small supply of wealth in money, especially in hyperinflation situations where money devalues quickly.[12]
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+ The limitations of barter are often explained in terms of its inefficiencies in facilitating exchange in comparison to money.
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+ It is said that barter is 'inefficient' because:
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+ Other anthropologists have questioned whether barter is typically between "total" strangers, a form of barter known as "silent trade". Silent trade, also called silent barter, dumb barter ("dumb" here used in its old meaning of "mute"), or depot trade, is a method by which traders who cannot speak each other's language can trade without talking. However, Benjamin Orlove has shown that while barter occurs through "silent trade" (between strangers), it also occurs in commercial markets as well. "Because barter is a difficult way of conducting trade, it will occur only where there are strong institutional constraints on the use of money or where the barter symbolically denotes a special social relationship and is used in well-defined conditions. To sum up, multipurpose money in markets is like lubrication for machines - necessary for the most efficient function, but not necessary for the existence of the market itself."[14]
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+
22
+ In his analysis of barter between coastal and inland villages in the Trobriand Islands, Keith Hart highlighted the difference between highly ceremonial gift exchange between community leaders, and the barter that occurs between individual households. The haggling that takes place between strangers is possible because of the larger temporary political order established by the gift exchanges of leaders. From this he concludes that barter is "an atomized interaction predicated upon the presence of society" (i.e. that social order established by gift exchange), and not typical between complete strangers.[15]
23
+
24
+ As Orlove noted, barter may occur in commercial economies, usually during periods of monetary crisis. During such a crisis, currency may be in short supply, or highly devalued through hyperinflation. In such cases, money ceases to be the universal medium of exchange or standard of value. Money may be in such short supply that it becomes an item of barter itself rather than the means of exchange. Barter may also occur when people cannot afford to keep money (as when hyperinflation quickly devalues it).[16]
25
+
26
+ An example of this would be during the Crisis in Bolivarian Venezuela, when Venezuelans resorted to bartering as a result of hyperinflation.[17]
27
+
28
+ Economic historian Karl Polanyi has argued that where barter is widespread, and cash supplies limited, barter is aided by the use of credit, brokerage, and money as a unit of account (i.e. used to price items). All of these strategies are found in ancient economies including Ptolemaic Egypt. They are also the basis for more recent barter exchange systems.[18]
29
+
30
+ While one-to-one bartering is practiced between individuals and businesses on an informal basis, organized barter exchanges have developed to conduct third party bartering which helps overcome some of the limitations of barter. A barter exchange operates as a broker and bank in which each participating member has an account that is debited when purchases are made, and credited when sales are made.
31
+
32
+ Modern barter and trade has evolved considerably to become an effective method of increasing sales, conserving cash, moving inventory, and making use of excess production capacity for businesses around the world. Businesses in a barter earn trade credits (instead of cash) that are deposited into their account. They then have the ability to purchase goods and services from other members utilizing their trade credits – they are not obligated to purchase from those whom they sold to, and vice versa. The exchange plays an important role because they provide the record-keeping, brokering expertise and monthly statements to each member. Commercial exchanges make money by charging a commission on each transaction either all on the buy side, all on the sell side, or a combination of both. Transaction fees typically run between 8 and 15%.[citation needed]
33
+
34
+ Throughout the 18th century, retailers began to abandon the prevailing system of bartering. Retailers operating out of the Palais complex in Paris, France were among the first in Europe to abandon the bartering, and adopt fixed-prices thereby sparing their clientele the hassle of bartering. The Palais retailers stocked luxury goods that appealed to the wealthy elite and upper middle classes. Stores were fitted with long glass exterior windows which allowed the emerging middle-classes to window shop and indulge in fantasies, even when they may not have been able to afford the high retail prices. Thus, the Palais-Royal became one of the first examples of a new style of shopping arcade, which adopted the trappings of a sophisticated, modern shopping complex and also changed pricing structures, for both the aristocracy and the middle classes.[19]
35
+
36
+ The Owenite socialists in Britain and the United States in the 1830s were the first to attempt to organize barter exchanges. Owenism developed a "theory of equitable exchange" as a critique of the exploitative wage relationship between capitalist and labourer, by which all profit accrued to the capitalist. To counteract the uneven playing field between employers and employed, they proposed "schemes of labour notes based on labour time, thus institutionalizing Owen's demand that human labour, not money, be made the standard of value."[20] This alternate currency eliminated price variability between markets, as well as the role of merchants who bought low and sold high. The system arose in a period where paper currency was an innovation. Paper currency was an IOU circulated by a bank (a promise to pay, not a payment in itself). Both merchants and an unstable paper currency created difficulties for direct producers.
37
+
38
+ An alternate currency, denominated in labour time, would prevent profit taking by middlemen; all goods exchanged would be priced only in terms of the amount of labour that went into them as expressed in the maxim 'Cost the limit of price'. It became the basis of exchanges in London, and in America, where the idea was implemented at the New Harmony communal settlement by Josiah Warren in 1826, and in his Cincinnati 'Time store' in 1827. Warren ideas were adopted by other Owenites and currency reformers, even though the labour exchanges were relatively short lived.[21]
39
+
40
+ In England, about 30 to 40 cooperative societies sent their surplus goods to an "exchange bazaar" for direct barter in London, which later adopted a similar labour note. The British Association for Promoting Cooperative Knowledge established an "equitable labour exchange" in 1830. This was expanded as the National Equitable Labour Exchange in 1832 on Grays Inn Road in London.[22] These efforts became the basis of the British cooperative movement of the 1840s. In 1848, the socialist and first self-designated anarchist Pierre-Joseph Proudhon postulated a system of time chits. In 1875, Karl Marx wrote of "Labor Certificates" (Arbeitszertifikaten) in his Critique of the Gotha Program of a "certificate from society that [the labourer] has furnished such and such an amount of labour", which can be used to draw "from the social stock of means of consumption as much as costs the same amount of labour."[23]
41
+
42
+ Michael Linton this originated the term "local exchange trading system" (LETS) in 1983 and for a time ran the Comox Valley LETSystems in Courtenay, British Columbia.[24] LETS networks use interest-free local credit so direct swaps do not need to be made. For instance, a member may earn credit by doing childcare for one person and spend it later on carpentry with another person in the same network. In LETS, unlike other local currencies, no scrip is issued, but rather transactions are recorded in a central location open to all members. As credit is issued by the network members, for the benefit of the members themselves, LETS are considered mutual credit systems.
43
+
44
+ The first exchange system was the Swiss WIR Bank. It was founded in 1934 as a result of currency shortages after the stock market crash of 1929. "WIR" is both an abbreviation of Wirtschaftsring (economic circle) and the word for "we" in German, reminding participants that the economic circle is also a community.[25]
45
+
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+ In Australia and New Zealand, the largest barter exchange is Bartercard, founded in 1991, with offices in the United Kingdom, United States, Cyprus, UAE and Thailand.[26] Other than its name suggests, it uses an electronic local currency, the trade dollar.
47
+
48
+ In business, barter has the benefit that one gets to know each other, one discourages investments for rent (which is inefficient) and one can impose trade sanctions on dishonest partners.[27]
49
+
50
+ According to the International Reciprocal Trade Association, the industry trade body, more than 450,000 businesses transacted $10 billion globally in 2008 – and officials expect trade volume to grow by 15% in 2009.[28]
51
+
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+ It is estimated that over 450,000 businesses in the United States were involved in barter exchange activities in 2010. There are approximately 400 commercial and corporate barter companies serving all parts of the world. There are many opportunities for entrepreneurs to start a barter exchange. Several major cities in the U.S. and Canada do not currently have a local barter exchange. There are two industry groups in the United States, the National Association of Trade Exchanges (NATE) and the International Reciprocal Trade Association (IRTA). Both offer training and promote high ethical standards among their members. Moreover, each has created its own currency through which its member barter companies can trade. NATE's currency is known as the BANC and IRTA's currency is called Universal Currency (UC).[29]
53
+
54
+ In Canada, barter continues to thrive. The largest b2b barter exchange is Tradebank, founded in 1987. P2P bartering has seen a renaissance in major Canadian cities through Bunz - built as a network of Facebook groups that went on to become a stand-alone bartering based app in January 2016. Within the first year, Bunz accumulated over 75,000 users[30] in over 200 cities worldwide.
55
+
56
+ Corporate barter focuses on larger transactions, which is different from a traditional, retail oriented barter exchange. Corporate barter exchanges typically use media and advertising as leverage for their larger transactions. It entails the use of a currency unit called a "trade-credit". The trade-credit must not only be known and guaranteed but also be valued in an amount the media and advertising could have been purchased for had the "client" bought it themselves (contract to eliminate ambiguity and risk).[citation needed]
57
+
58
+ Soviet bilateral trade is occasionally called "barter trade", because although the purchases were denominated in U.S. dollars, the transactions were credited to an international clearing account, avoiding the use of hard cash.
59
+
60
+ In the United States, Karl Hess used bartering to make it harder for the IRS to seize his wages and as a form of tax resistance. Hess explained how he turned to barter in an op-ed for The New York Times in 1975.[31] However the IRS now requires barter exchanges to be reported as per the Tax Equity and Fiscal Responsibility Act of 1982. Barter exchanges are considered taxable revenue by the IRS and must be reported on a 1099-B form. According to the IRS, "The fair market value of goods and services exchanged must be included in the income of both parties."[32]
61
+
62
+ Other countries, though, do not have the reporting requirement that the U.S. does concerning proceeds from barter transactions, but taxation is handled the same way as a cash transaction. If one barters for a profit, one pays the appropriate tax; if one generates a loss in the transaction, they have a loss. Bartering for business is also taxed accordingly as business income or business expense. Many barter exchanges require that one register as a business.
63
+
64
+ In Spain (particularly the Catalonia region) there is a growing number of exchange markets.[33] These barter markets or swap meets work without money. Participants bring things they do not need and exchange them for the unwanted goods of another participant. Swapping among three parties often helps satisfy tastes when trying to get around the rule that money is not allowed.[34]
65
+
66
+ Other examples are El Cambalache in San Cristobal de las Casas, Chiapas, Mexico[35] and post-Soviet societies.[36]
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1
+ Troy (Ancient Greek: Τροία, Troía, Ἴλιον, Ílion or Ἴλιος, Ílios; Latin: Troia and Ilium;[note 1] Hittite: 𒌷𒃾𒇻𒊭 Wilusa or 𒋫𒊒𒄿𒊭 Truwisa;[3][4] Turkish: Truva or Troya) was a city in the northwest of Asia Minor (modern Turkey), southwest of the Canakkale Strait, south of the mouth of the Dardanelles and northwest of Mount Ida.[note 2] The location in the present day is the hill of Hisarlik and its immediate vicinity. In modern scholarly nomenclature, the Ridge of Troy (including Hisarlik) borders the Plain of Troy, flat agricultural land, which conducts the lower Scamander River to the strait. Troy was the setting of the Trojan War described in the Greek Epic Cycle, in particular in the Iliad, one of the two epic poems attributed to Homer. Metrical evidence from the Iliad and the Odyssey suggests that the name Ἴλιον (Ilion) formerly began with a digamma: Ϝίλιον (Wilion);[note 3] this is also supported by the Hittite name for what is thought to be the same city, Wilusa.
2
+
3
+ After a destruction at the end of the Bronze Age, believed to represent the end of the Trojan War, and a period of abandonment or near-abandonment during the subsequent Dark Age, the site acquired a new population of Greek-speakers, who built a classical city that became along with the rest of Anatolia a part of the Persian Empire. The Troad was liberated by Alexander the Great, an admirer of Achilles, who he believed had the same type of glorious (but short-lived) destiny. After the Roman conquest of this now Hellenistic Greek-speaking world, a new capital called Ilium (from Greek: Ἴλιον, Ilion) was founded on the site in the reign of the Roman Emperor Augustus. It flourished until the establishment of Constantinople, became a bishopric, was abandoned, repopulated for a few centuries in the Byzantine era, was abandoned again, and is now a Latin Catholic titular see. Most recently it has risen to prominence as an archaeological site.
4
+
5
+ In the mid-19th century the Calvert family, wealthy Levantine English settlers of the Troad, occupying a working farm a few miles from Hisarlik, purchased much of the hill in the belief that it contained the ruins of Troy. They were antiquarians. Two of the family, Frederick and especially the youngest, Frank, surveyed the Troad and conducted a number of trial excavations there. In 1865, Frank Calvert excavated trial trenches on the hill, discovering the Roman settlement. Realizing he did not have the funds for a full excavation, he attempted to recruit the British Museum, and was refused. A chance meeting with Calvert in Çanakkale and a visit to the site by Heinrich Schliemann, a wealthy German businessman and archaeologist, also looking for Troy, offered a second opportunity for funding.[5] Schliemann had been at first skeptical about the identification of Hisarlik with Troy, but was persuaded by Calvert.[6] As Schliemann was about to leave the area, Calvert wrote to him asking him to take over the entire excavation. Schliemann agreed. The Calverts, who made their money in the diplomatic service, expedited the acquisition of a Turkish firman. In 1868, Schliemann excavated an initial deep trench across the mound called today "Schliemann's trench." These excavations revealed several cities built in succession. Subsequent excavations by following archaeologists elaborated on the number and dates of the cities.
6
+
7
+ Today a small village near the ruins, Tevfikiye, supports the tourist trade visiting the Troia archaeological site. It lies within the province of Çanakkale, some 30 kilometres (19 mi) south-west of the provincial capital, also called Çanakkale. The current map shows Ilium a little way inland from the Scamander estuary across the plain of Troy. According to Korfmann, due to Troy's location near the Aegean Sea, the Sea of Marmara, and the Black Sea, it was a central hub for military activities and trade, and the chief site of a culture he calls the "Maritime Troja Culture," which extended over the region between the Black and Aegean Seas.[7]
8
+
9
+ Troy was added to the UNESCO World Heritage list in 1998.
10
+
11
+ Homeric Troy refers primarily to the city described in the Iliad, the earliest literary work in Europe[citation needed]. This is a long originally oral poem in its own dialect of ancient Greek in dactylic hexameter, in tradition composed by a blind poet of the Anatolian Greek coast, Homer. It covers the 10th year of a war against Troy conducted by a coalition of Achaean, or Greek, states under the leadership of a high king, Agamemnon of Mycenae. The city was defended by a coalition of states in the Dardanelles and West Anatolian region under another high king, Priam, whose capital was Troy. The cause of the war was the elopement of Agamemnon's brother's wife, Helen, with Paris, a prince of Troy.
12
+
13
+ After the literary time of the poem, the city was destroyed when the Greeks pretended to leave after secreting a squad of soldiers in a gigantic wooden horse monument, which the Trojans brought inside the walls.[note 4] In the dead of night they exited the horse and opened the gates to the Achaeans nearby. Troy was burned and the population slaughtered, although many had other fates.
14
+
15
+ Besides the Iliad, there are references to Troy in the other major work attributed to Homer, the Odyssey, as well as in other ancient Greek literature (such as Aeschylus's Oresteia). The Homeric legend of Troy was elaborated by the Roman poet Virgil in his Aeneid. The fall of Troy with the story of the Trojan Horse and the sacrifice of Polyxena, Priam's youngest daughter, is the subject of a later Greek epic by Quintus Smyrnaeus ("Quintus of Smyrna").
16
+
17
+ The Greeks and Romans took for a fact the historicity of the Trojan War and the identity of Homeric Troy with a site in Anatolia on a peninsula called the Troad (Biga Peninsula). Alexander the Great, for example, visited the site in 334 BC and there made sacrifices at tombs associated with the Homeric heroes Achilles and Patroclus. In Piri Reis book Kitab-ı Bahriye (Book of the Sea, 1521) which details many ports and islands of the Mediterranean, the description of the island called Tenedos mentions Troy and its ruins, lying on the shore opposite of the island.[8]
18
+
19
+ The fact that the topography around the site matches the topographic detail of the poem gives to the poem a life-like quality not equaled by other epics. In the Iliad, the Achaeans set up their camp near the mouth of the River Scamander (modern Karamenderes),[9] where they beached their ships. The city of Troy itself stood on a hill, across the plain of Scamander, where the battles of the Trojan War took place. The site of the ancient city is some 5 kilometres (3.1 mi) from the coast today, but 3,000 years ago the mouths of Scamander were much closer to the city,[10] discharging into a large bay that formed a natural harbor, which has since been filled with alluvial material. Recent geological findings have permitted the identification of the ancient Trojan coastline, and the results largely confirm the accuracy of the Homeric geography of Troy.[11]
20
+
21
+ In November 2001, the geologist John C. Kraft from the University of Delaware and the classicist John V. Luce from Trinity College, Dublin, presented the results of investigations, begun in 1977, into the geology of the region.[12] They compared the present geology with the landscapes and coastal features described in the Iliad and other classical sources, notably Strabo's Geographia, and concluded that there is a regular consistency between the location of Schliemann's Troy and other locations such as the Greek camp, the geological evidence, descriptions of the topography and accounts of the battle in the Iliad.[13]
22
+
23
+ The Dark Age following the fall of Troy is called so because for a time writing in Greece disappeared. There are consequently no historians from the period. Writing reappeared in the Archaic Period, after which, in the Classical Period, many historians turned their pens to record such histories of the Trojan War as had survived in oral tradition. They offer a span of about two centuries from the 1334 BC date of Duris of Samos to the 1135 BC date of Ephoros of Kyme in Aeolis. Blegen preferred the 1184 BC date of Eratosthenes, which was in his day the most favored.[14][note 5] Whether or not the archaeology matched this span and these dates was to be determined by excavation.
24
+
25
+ With the rise of critical history, Troy and the Trojan War were consigned to legend.[note 6] However, not everyone agreed with this view. The dissidents were to become the first archaeologists at Troy. The true location of ancient Troy had for centuries remained the subject of interest and speculation.[15] Travellers in Anatolia looked for possible locations. Because of its name, the Troad peninsula was highly suspect.
26
+
27
+ Early modern travellers in the 16th and 17th centuries, including Pierre Belon and Pietro Della Valle, had identified Troy with Alexandria Troas, a ruined town approximately 20 kilometres (12 mi) south of the currently accepted location.[16] In the late 18th century, Jean Baptiste LeChevalier identified a location near the village of Pınarbaşı, Ezine, a mound approximately 5 kilometres (3.1 mi) south of the currently accepted location. Published in his Voyage de la Troade, it was the most commonly accepted theory for almost a century.[17]
28
+
29
+ In 1822, the Scottish journalist Charles Maclaren was the first to identify with confidence the position of the city as it is now known.[18][19] In the second half of the 19th century archaeological excavation of the site believed to have been Homeric Troy began. As the Iliad is taught in every Greek language curriculum in the world, interest in the site has been unflagging. Homeric experts often memorize large parts of the poem. Literary quotes are commonplace. Since the Calvert family began excavation at Hisarlik, hundreds of interested persons have excavated there. Fortunately all excavation has been conducted under the management of key persons termed its "archaeologists." Their courses of excavation have been divided into the phases described below. Sometimes there have been decades between phases. Today interest in the site is as strong as ever. Further plans for excavation have no end in the foreseeable future.
30
+
31
+ Frank Calvert was born into an English Levantine family on Malta in 1828. He was the youngest of six sons and one daughter born to James Calvert and his wife, the former Louisa Lander, the sister of Charles Alexander Lander, James' business partner. In social standing they were of the aristocracy. James was a distant relative of the Calverts who founded Baltimore, Maryland,[21] and Louisa was a direct descendant of the Campbells of Argyll (Scottish clansmen).[22] Not having inherited any wealth, they took to the colonies, married in Ottoman Smyrna in 1815, and settled in Malta, which had changed hands from the French to the British Empire with the Treaty of Paris (1814). They associated with the "privileged" social circles of Malta, but they were poor. James clerked in the mail and grain offices of the Civil Service.[23]
32
+
33
+ The family regarded itself as a single enterprise. They shared property, assisted each other, lived together and had common interests, one of which was the antiquities of the Troad. They did not do well in Malta, but in 1829 the Dardanelles region underwent an upswing of its business cycle due to historical circumstances. The Greek War of Independence was about to be concluded in favor of an independent state by the Treaty of Constantinople (1832). The Levant Company, which had had a monopoly on trade through the Dardanelles, was terminated. The price in pounds of the Turkish piastre fell. A manyfold increase in British traffic through straits was anticipated. A new type of job suddenly appeared: British Consul in the Dardanelles, which brought wealth with it.[24][note 7]
34
+
35
+ Charles Lander applied, and was made British Consul of the Dardanelles in 1829. He spoke five languages, knew the region well, and had the best connections. A row of new consular offices was being constructed in Çanakkale along the shore of the strait. He was at first poor. In 1833 he bought a house in town ample enough to invite his sister's sons to join him in the enterprise. Without exception they left home at 16 to be tutored in the trade at their uncle's house and placed in lucrative consular positions. Frederick, the eldest, stayed on to assist Charles. The youngest, Frank, at school in Athens, arrived last, but his interest in archaeology led him into a different career.[21]
36
+
37
+ Çanakkale was a boom town. In 1831 Lander married Adele, a brief but idyllic relationship that gave them three daughters in quick succession. When the Calverts began to arrive, finding quarters in the crowded town proved to be difficult. The Turkish building code requiring buildings of wood, conflagrations were frequent.[25] The family escaped one fire with nothing but the clothes they were wearing.[26] Lander's collection of books on the Troad was totally destroyed. In 1840 Lander suffered a tragedy when his wife, Adele, died in her 40's, leaving three small children. He chose this time to settle his estate, making Frederick his legal heir, guardian of his children, and co-executor (along with himself).
38
+
39
+ Lander dedicated himself to the consular service, leaving the details of the estate and its reponsibilites to Frederick. The family grew wealthy on the fees paid by the ships they serviced. When Frank arrived in 1845[27] with his sister he had nothing much to do. By this time the family had a new library. Using its books Frank explored the Troad.[28] He and Lander became collectors. The women in the family took a supportive role as well.
40
+
41
+ Lander died in 1846 of a fever endemic to the region, leaving Frederick as executor of the will and head of the family. In 1847 he assumed his uncle's consular position. He was also an agent of Lloyd's of London, which insured ship cargos. Despite Frank's youth he began to play an important role in the family consular business, especially when Frederick was away.[29] A few years prior to the death of Lander, the population of Çanakkale was on the rise, from 10,000 in 1800 to 11,000 in 1842.[30] The British numbered about 40 families.[31] The increase in ship traffic meant prosperity for the Calverts, who expedited the ships of several nations, including the United States. They had other ambitions: James William Whittall, British consul in Smyrna, was spreading his doctrine of the "Trojan Colonization Society," (never more than an idea) which was influential on the Calverts, whom he visited.[32]
42
+
43
+ In 1847 Frederick invested the profits of the family business in two large tracts in the Troad, amounting to many thousands of acres.[33][note 8] He founded a company, Calvert Bros. and Co., an "extended family company."[34] The first purchase was a farm at Erenkőy, on the coast about half-way between Çanakkale and Troy. Frederick used it as a station for ships that could not make Çanakkale. The area was a target for Greek immigration. The family became money-lenders, lending only to Greeks at rates considered high (20%).[35]
44
+
45
+ Frederick also bought a farm he intended to work, the Batak Farm (named for the Batak wetlands), later changed by Frank to Thymbra Farm, because he believed it was the site of Homeric Thymbra, after which the Thymbra Gate of Homeric Troy had been named. It was located at an abandoned village called Akça Köy, 4 mi. to the southeast of Hisarlik. The farm was the last of the village. It harvested and marketed the cups and acorns of Quercus macrolepis, the Valonia Oak, from which valonia, a compound used in dying and tanning, is extracted. The farm also raised cotton and wheat and bred horses. Frederick introduced the English plough and drained the wetlands. The farm eventually became famous as a way station for archaeologists and the home of the Calvert collection of antiquities, which Frank kept locked in a hidden room. The main house, featuring multiple guest bedrooms, was situated on a low ridge in a compound with several outbuildings. It was more of a manor, operated by farm workers and domestic servants.
46
+
47
+ In 1850–1852 Frederick solved the residence problem by having a mansion built for him in Çanakkale. Two Turkish houses were said to have been put together, but Turkish houses were required to be of wood. This one was of massive stone, which was permitted to foreigners, and was placed partly on fill jetting into the straits. It probably was the length of two Turkish houses. It remained the major building of the town until it was removed in 1942, due to earlier earthquake damage. The last of the Calvert descendants still in the region had ceded it to the town in 1939. The Town Hall was then built on the site. The mansion's extensive gardens became a public park.[37]
48
+
49
+ The entire family of the times took up permanent residence in the mansion, which was never finished. It was almost always occupied by visitors and social events. The Calverts began a tour-guide business, conducting visitors throughout the Troad. Frank was the chief guide. The women held musicales and sang in the salons. The house attracted a stream of distinguished visitors, each with a theory about the location of Troy. Frederick, however, was not there for the opening of the house. After a fall from a horse in 1851, complications forced him to seek medical care in London for 18 months,[38] the first of a series of disasters. He was back by 1853.
50
+
51
+ The Crimean War began in October 1853 and lasted through February 1856. Russia had arbitrarily occupied the Danube frontier of the Ottoman Empire including the Crimea, and Britain and France were providing military assistance to the Ottomans. The rear of the conflict was Istanbul and the Dardanelles. Britain relied heavily on the Levantine families for interfacing, intelligence, and guidance. Edmund Calvert was a British agent, but this was not Frederick's calling. Not long after his return the initial British expeditionary force of 10,000 men was held up in ships in the straits, with no place to bivouac, no supplies, and a commissariat of four non-Turkish speakers.[39]
52
+
53
+ The British Army had reached a low point of efficiency since Wellington.[40] Although it was the responsibility of Parliament, the fact that the crown retained the prerogative of command made them hesitate to update it, for fear of its being used against them.[41] One of the major problems was the fragmentation of the administration into "a number of separate, distinct, and mutually independent authorities," with little centralization.[42] There were always issues of who was in command and what they commanded. A Supply Corps as such did not exist. The immediate needs of the soldiers were supplied by the Commissariat Department, responsible to the Treasury.[43] Commissaries were assigned to units as needed, but they acted to solve supply problems ad hoc. They had no idea beforehand what the army needed, or what it had, or where it was located.
54
+
55
+ All the needs were given to contractors, who usually required money in advance. They were allowed to borrow from recommended banks. The Commissariat then paid the banks, but should it fail to do so, the debts were still incumbent on the debtors. Contractors were allowed to charge a percentage for their services, and also to include a percentage given to their suppliers as enticement. The Commissariat could thus build entire impromptu supply departments on the basis of immediate need, which is what Frederick did for them.[44]
56
+
57
+ The logistics problems were of the same type customarily undertaken by the consular staff, but of larger scale. Frederick was able to perform critical services for the army. Within several days he had all the men billeted ashore and had developed an organization of local suppliers on short notice. He secured their immediate attention by offering higher interest rates, to which the Commissary did not then object. He was so successful that he was given the problem of transporting men and supplies to the front.[note 9] For that he developed his own transport division of contractors paid as direct employees of his own company. He also advised the Medical Department in their choice of a site near Erenköy for a military hospital, named Renkioi Hospital.[44]
58
+
59
+ The army, arriving at Gallipoli in April, 1854, did well at first, thanks to the efforts of Frederick Calvert and his peers. They were contracted by Deputy Assistant Commander-General of the Commissariat, John William Smith, on the instruction of the Commander-General, William Filder, who had given Smith their names in advance, especially that of Frederick Calvert. Frederick was waiting for the fleet in Gallipoli.[45][note 10] By June the army was doing badly. The Commissary seemed to have no understanding of military schedules. Needed supplies were not getting to their destinations for a number of reasons: perishables were spoiled through delay, cargos were lost or abandoned because there was no tracking system, or cut because a commissary speculated that they should be, etc. Frederick attempted to carry on by using his own resources in the expectation of collecting the money later by due process. By the end of the war his bill to the Commissary would be several thousand pounds. He had had to mortgage family properties in the Troad.[46]
60
+
61
+ By June it was obvious to Parliament that the cabinet position of “Secretary of State for War and the Colonies” was beyond the ability of only one minister. He was divested of his colonial duties, leaving him as Secretary of State for War,[47] but the Commissary was still not in his domain. In August, Frederick purchased the winter feed for the animals and left it on the dock at Salonica. Filder had adopted a policy of purchasing hay from London and having it pressed for land transport, even though chopped hay was readily available at a much cheaper price around the Dardanelles.[48] The Commissariat was supposed to inspect and accept it at Salonica, but the presses had been set up in the wrong location. By the time they were ready for the hay, most of it had spoiled, so they did not accept any of it.
62
+
63
+ The winter was especially severe. The animals starved, and without transport, so did the men, trying to make do without food, clothing, shelter or medical supplies.[49] Estimates of the death rate were as high as 35%, 42% in the field hospitals.[50] Florence Nightingale on the scene sounded the alarm to the general public. A scandal ensued; Prince Albert wrote to the Prime Minister. The folly of an army dying because not allowed to help itself while its Commissariat was not efficient enough to move even the minimum of supplies became manifest to the whole nation. In December Parliament placed the Commissariat under the army and opened an investigation.[51] In January, 1855, the government resigned, to be replaced shortly by another determined to do whatever was necessary to obtain a functional supply corps.[52]
64
+
65
+ The army found that it could not after all dispense with the Treasury or its system of payment. The first investigation went before Parliament in April, 1855. Filder’s defense was that he had conformed strictly to regulations,[note 11] and that he was not responsible for accidental events, which were “the visitations of God.”[53] John William Smith, Frederick’s handler in the Commissariat, included a number of favorable statements about him in the report, such as “the Commissariat would have been perfectly helpless without Mr. Calvert.”[54] Parliament exonerated the Commissariat, finding “no one in the Crimea was to blame.”[55]
66
+
67
+ Anticipating this result, the new government started a secret investigation of its own under J. McNeill, a civilian physician, and a milItary officer, Colonel A.M. Tulloch, which it outed in April after the acquittal. The new investigation lasted until January, 1856, and had nothing favorable to say. Losses higher than any battle could produce, and higher than those of any of the allies, were not to be dismissed as accidental.
68
+
69
+ The new commissioners attacked the system: “the system hitherto relied on as sufficient to provide for every emergency, had totally failed.”[56] The blow fell mainly on Filder. He had plenty of alternatives, Tulloch asserted, which he might have been expected to take. Chopped hay and cattle were readily and cheaply available in the Constantinople region. Filder had some cattle transports at his command in October. Once the supplies had been transported to the Crimea, they could have been carried inland by the troops themselves.[57] Of Filder, Tulloch said: “He was highly paid — not to do merely what he was ordered, but in the expectation that, when difficulties arose, he would show himself equal to the emergency, by ... exercising that discretion and intelligence which the public has a right to expect ....”[58]
70
+
71
+ Filder was retired by the medical board because of age and sent home. Meanwhile the Commissary had introduced the word "profiteering" in a effort to cast the blame from itself. The decisions had been made by greedy contractors charging high interest rates, who had introduced delays to push the price up. John William Smith recanted what he had said about Frederick, now claiming that Frederick had put private interests before the public, without clarifying what he meant. The insinuation was enough to brand him as a profiteer.[59]. The entire Commissariat took it up as a theme, the banks refusing to honor contractor claims. Restrictions on loans tightened; cash flow problems developed. The inflated economy of the Troad began to collapse. The report was released in January. By then most contractors were in bankruptcy. British troops went home at the end of the war in February, having turned the Turkish merchants in the Troad against the English.
72
+
73
+ The cost of living remained high. Frederick was no longer trusted as a consular agent and had trouble finding work. His friend, John Brunton, head of the military hospital near Erenköy, was ordered to dismantle and sell the facility. He suggested that Brunton sell the medical supplies to him as surplus at a discount, so that he could recoup some of his estate by reselling them. Turning on him, as Smith had done, Brunton denounced him publicly.
74
+
75
+ Criminal charges were brought against Frederick for non-payment of debt to the War Office by the Supreme Consular Court of Istanbul in March, 1857. Due to difficulty in proving their case, it went on for months, being finally transferred to London,[60] where Frederick joined it in February, 1858. In 1859 he served a prison term of ten weeks on one debt. Subsequently the Foreign Office stepped in to manage his appeal. The military had not understood how the interest system worked. He won his case before Parliament, with commendation and thanks, and payment of the several thousand plus backpay and interest, arriving home 2.5 years after he had left it, to rescue the estate.[note 12]
76
+
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+ During the 1860s Frederick Calvert's life and career were mainly consumed by a case of insurance fraud termed by the press the “Possidhon affair." An attempt was made to defraud Lloyd's of London of payments to an imaginary person claiming to own an imaginary ship, the Possidhon, that had gone to the bottom when its imaginary cargo burned, a claim made through Frederick. The perpetrators of the fraud, originally the witnesses of the fire, named Frederick as their ringleader. The trial was not a proper one, and Frederick was convicted on technicalities. He protested that he was the victim of an Ottoman frame-up, and was supported in that plea by his brother, Frank. There were a number of circumstances that remain historically unexplained. Modern historians who think he was guilty characterize him as a charismatic profiteer of shady ethics, while those who think he was innocent point to his patriotic motives in helping the British Army to the detriment of his own estate and his acquittal by Parliament.
78
+
79
+ Having returned from London in October, 1860, with enough money to restore the family estate, Frederick now turned his attention to the family avocation, archaeology, rejecting a lucrative job offer as a Consul in Syria.[61] Frank, now age 32, had long been the master of the estate and of the business. By this time he was also a skilled and respected archaeologist. He spent all of his spare time investigating and excavating the numerous habitation and burial sites of the Troad. He was an invaluable consultant to specialists in many areas from plants to coins. Frederick joined him in this life by choice. For a few years he was able to work with Frank in expanding Lander’s library and collection, and in exploring and excavating ancient sites.
80
+
81
+ In 1846 Frederick married Eveline, an heiress of the wealthy Abbotts, owners of some mines in Turkey. They had at least five known children.
82
+
83
+ Frederick’s wife’s uncle, William Abbott, had gone with him to London, where they purchased a house for mutual residence. Frederick set him up in a few different businesses, the last being Abbott Brothers, dealers in firewood. His son, however, William George Abbott, a junior partner of Frederick in the consular business, remained in the Dardanelles to handle business there as acting consul.[note 13] In January, 1861, the consular office was approached by a Turkish merchant, Hussein Aga, requesting 12000 £. ($16000} of insurance from Lloyd’s on the cargo of the Possidhon, which was olive oil. He claimed to be a broker marketing the oil produced by certain pashas and now wished to sell it in Britain.[citation needed]
84
+
85
+ Frederick requested William in London to borrow money as Abbot Brothers to finance the premiums.[62] The debt was to be paid when the cargo was sold. It isn’t clear whether Abbott was to sell it, and if so, in whose name. The cargo, being insured by him, was consigned to him. A loan of 1500 £ ($2000) was effected on April 11, and the premiums were paid.
86
+
87
+ The ship, cleared to sail from Edremit to Britain by Frederick’s office on April 4, sailed on the 6th. Frederick was to have inspected it before issuing the clearance, but he did not. On April 28 Frederick notified Lloyd’s by telegram that the vessel had been seen burning off Lemnos in a heavy wind on April 8, which is peculiar, because it ought to have been far from Lemnos by then. When it had not arrived months later the creditors for the premiums requested their money. Frederick submitted a claim through Abbott for a total loss. He suggested Greek pirates and collaboration of the crew as causes, implicating Hussein Aga, who had not been seen since then. Lloyd’s requested documents giving testimony of the loss, turning the case over to Lloyd's Salvage Association.
88
+
89
+ Frederick forwarded to Abbott in London four affadavits from British consular agents on Tenedos and Samos of visual sightings of the ship. Conspicuously absent were any Turkish documents that should have been examined before permission to sail was granted. An investigator from Lloyd’s Salvage working from Constantinople finding no record of either Aga or the ship concluded to a fraud. Simultaneouly Frederick, conducting his own investigation, reached a similar conclusion. He had been duped by a person pretending to be a fictional Hussein Aga. The witnesses produced a confession, naming Frederick as mastermind of the scheme. The Salvage Association turned the matter over to the Foreign Office. M. Tolmides, consular agent at Tenedos, admitted to signing the affadavits. His defense was that he had given Frederick blank signed forms.
90
+
91
+ The Foreign Office issued a public statement questioning Frederick's credibility. He requested permission to leave his post to travel to London to defend himself. Permission was denied. On April 30 he issued a statement that he had been set up and was being framed by an unknown agent, for whom he was conducting an unsuccessful search at Smyrna. He found some support in the British ambassador, Henry Bulwer, 1st Baron Dalling and Bulwer, a liberal and a freemason, who accepted him as credible, and noted the hostility of Turkish officialdom against him. However, unless Frederick could produce some evidence of the conspiracy, he affirmed, he would officially have to side with the insurance company. The matter became international. Turkish harbor officials claimed, via Lloyd’s agents, that Frederick had submitted forged documents to them. The Ottoman Porte compalined. The Prince of Wales scheduled a visit. Fredrick was going to be brought before a consular court, an agency with a reputation for corruption; in particular, bribability.[citation needed]
92
+
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+ Due to the publicity skills of Heinrich Schliemann and the public discreditation of Frederick as a convicted felon, the contributions mainly of Frank to the excavation of Troy remained unknown and unappreciated until the end of the 20th century, when the Calverts became an object of special study. A number of misunderstandings still cling to them. One is that Schliemann discovered Troy on land he had the foresight to purchase from the Calverts. To the contrary, it was Frank who convinced Frederick to purchase Hissarlik as the probable site of Troy, and Frank who convinced Schliemann that it was there, and to partner with him in its excavation.[63] The Calverts did not hand anything over; they remained on site excavating with him and attempting to advise and manage him. Frank was often a sharp critic. Frank is sometimes called "self-taught." Educationally this was not true. He did not attend university, but there would have been no point, as archaeology was not yet taught there. Frank was the first modern (19th century) to excavate in the Troad.[64] He knew more than all the visitors he tutored.
94
+
95
+ In 1866, Frank Calvert, the brother of the United States' consular agent in the region, made extensive surveys and published in scholarly journals his identification of the hill of New Ilium (which was on farmland owned by his family) on the same site. The hill, near the city of Çanakkale, was known as Hisarlik.[65]
96
+
97
+ The British diplomat, considered a pioneer for the contributions he made to the archaeology of Troy, spent more than 60 years in the Troad (modern day Biga peninsula, Turkey) conducting field work.[66] As Calvert was a principal authority on field archaeology in the region, his findings supplied evidence that Homeric Troy might have existed on the hill, and played a major role in convincing Heinrich Schliemann to dig at Hisarlik.[20]
98
+
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+ In 1868, German archaeologist Heinrich Schliemann visited Calvert and secured permission to excavate Hisarlik. He sincerely believed that the literary events of the works of Homer could be verified archaeologically. A divorced man in his 40's who had acquired some wealth as a merchant in Russia, he decided to use the wealth to follow his boyhood interest in finding and verifying the city of Troy. Leaving his former life behind, he advertised for a wife whose skills and interest were on a par with his own, Sophia. She was 17 at the time but together they excavated Troy, sparing no expense.
100
+
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+ Heinrich began by excavating a trench across the mound of Hisarlik to the depth of the settlements, today called "Schliemann's Trench." In 1871–73 and 1878–79, he discovered the ruins of a series of ancient cities dating from the Bronze Age to the Roman period. He declared one of these cities—at first Troy I, later Troy II—to be the city of Troy, and this identification was widely accepted at that time. Subsequent archaeologists at the site were to revise the date upward; nevertheless, the main identification of Troy as the city of the Iliad, and the scheme of the layers, have been kept.
102
+
103
+ Some of Schliemann's portable finds at Hisarlik have become known as Priam's Treasure, such as the jewelry photographed displayed on Sophia. The artifacts were acquired from him by the Berlin museums. As Sophia matured she became an invaluable assistant to Schliemann, whom he employed especially in social situations requiring the use of modern Greek. After his death she became caretaker of his funds and publications, continuing to advocate for his beliefs. She was a respected socialite in Athens.
104
+
105
+ Wilhelm Dörpfeld (1893–94) joined the excavation at the request of Schliemann. After Schliemann left, he inherited the management of it. His chief contribution was the detailing of Troy VI. He published his findings separately.[67]
106
+
107
+ Carl Blegen, professor at the University of Cincinnati, managed the site 1932–38. These archaeologists, though following Schliemann's lead, added a professional approach not available to Schliemann. He showed that there were at least nine cities. In his research, Blegen came to a conclusion that Troy's nine levels could be further divided into forty-six sublevels,[68] which he published in his main report.[69]
108
+
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+ In 1988, excavations were resumed by a team from the University of Tübingen and the University of Cincinnati under the direction of Professor Manfred Korfmann, with Professor Brian Rose overseeing Post-Bronze Age (Greek, Roman, Byzantine) excavation along the coast of the Aegean Sea at the Bay of Troy. Possible evidence of a battle was found in the form of bronze arrowheads and fire-damaged human remains buried in layers dated to the early 12th century BC. The question of Troy's status in the Bronze-Age world has been the subject of a sometimes acerbic debate between Korfmann and the Tübingen historian Frank Kolb in 2001–2002.
110
+
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+ Korfmann proposed that the location of the city (close to the Dardanelles) indicated a commercially oriented city that would have been at the center of a vibrant trade between the Black Sea, Aegean, Anatolian and Eastern Mediterranean regions. Kolb disputed this thesis, calling it "unfounded" in a 2004 paper. He argues that archaeological evidence shows that economic trade during the Late Bronze Age was quite limited in the Aegean region compared with later periods in antiquity. On the other hand, the Eastern Mediterranean economy was more active during this time, allowing for commercial cities to develop only in the Levant. Kolb also noted the lack of evidence for trade with the Hittite Empire.[70]
112
+
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+ In August 1993, following a magnetic imaging survey of the fields below the fort, a deep ditch was located and excavated among the ruins of a later Greek and Roman city. Remains found in the ditch were dated to the late Bronze Age, the alleged time of Homeric Troy. Among these remains are arrowheads and charred remains.[71] It is claimed by Korfmann that the ditch may have once marked the outer defenses of a much larger city than had previously been suspected. In the olive groves surrounding the citadel, there are portions of land that were difficult to plow, suggesting that there are undiscovered portions of the city lying there. The latter city has been dated by his team to about 1250 BC, and it has been also suggested—based on recent archeological evidence uncovered by Professor Manfred Korfmann's team—that this was indeed the Homeric city of Troy.
114
+
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+ Helmut Becker utilized magnetometry in the area surrounding Hisarlik. He was conducting an excavation in 1992 to locate outer walls of the ancient city. Becker used a caesium magnetometer. In his and his team's search, they discovered a "'burnt mudbrick wall' about 400 metres south of the Troy VI fortress wall."[72] After dating their find, it was deemed to have been from the late Bronze Age, which would put it either in Troy VI or early Troy VII. This discovery of an outer wall away from the tell proves that Troy could have housed many more inhabitants than Schliemann originally thought.
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+
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+ In summer 2006, the excavations continued under the direction of Korfmann's colleague Ernst Pernicka, with a new digging permit.[73]
118
+
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+ In 2013, an international team made up of cross-disciplinary experts led by William Aylward, an archaeologist at the University of Wisconsin-Madison, was to carry out new excavations. This activity was to be conducted under the auspices of Çanakkale Onsekiz Mart University and was to use the new technique of "molecular archaeology".[74] A few days before the Wisconsin team was to leave, Turkey cancelled about 100 excavation permits, including Wisconsin's.[75]
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+
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+ In March 2014, it was announced that a new excavation would take place to be sponsored by a private company and carried out by Çanakkale Onsekiz Mart University. This will be the first Turkish team to excavate and is planned as a 12-month excavation led by associate professor Rüstem Aslan. The University's rector stated that "Pieces unearthed in Troy will contribute to Çanakkale’s culture and tourism. Maybe it will become one of Turkey’s most important frequented historical places.”[76]
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+
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+ The Turkish government created the Historical National Park at Troy on September 30, 1996. It contains 136 square kilometres (53 sq mi) to include Troy and its vicinity, centered on Troy.[77] The purpose of the park is to protect the historical sites and monuments within it, as well as the ecology of the region. In 1998 the park was accepted as a UNESCO World Heritage Site.
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+
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+ In 2015 a Term Development Revision Plan was applied to the park. Its intent was to develop the park into a major tourist site.[78] Plans included marketing research to determine the features most of interest to the public, the training of park personnel in tourism management, and the construction of campsites and facilities for those making day trips. These latter were concentrated in the village of Tevfikiye, which shares Troy Ridge with Troy.
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+
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+ Public access to the ancient site is along the road from the vicinity of the museum in Tevfikiye to the east side of Hisarlik. Some parking is available. Typically visitors come by bus, which disembarks its passengers into a large plaza ornamented with flowers and trees and some objects from the excavation. In its square is a large wooden horse monument, with a ladder and internal chambers for use of the public. Bordering the square is the gate to the site. The public passes through turnstiles. Admission is usually not free. Within the site the visitors tour the features on dirt roads or for access to more precipitous features on railed boardwalks. There are many overlooks with multilingual boards explaining the feature. Most are outdoors, but a permanent canopy covers the site of an early megaron and wall.
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+
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+ The archaeological site of Troy was added to the UNESCO World Heritage list in 1998.
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+
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+ For a site to be named a UNESCO World Heritage Site, it must be claimed to have Outstanding Universal Value. This means that it must be historically, culturally, or scientifically significant to all peoples of the world in some manner. According to the UNESCO site on Troy, its historical significance was gained because the site displays some of the "first contact between...Anatolia and the Mediterranean world".[79] The site's cultural significance is gained from the multitudes of literature regarding the famed city and history over centuries. Many of the structures dating to the Bronze Age and the Roman and Greek periods are still standing at Hisarlik. These give archeological significance to the site as well.
132
+
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+ In 2018 the Troy Museum (Turkish Troya Müzesi) was opened at Tevfikiye village 800 metres (870 yd) east of the excavation. A design contest for the architecture had been won by Yalin Mimarlik in 2011. The cube-shaped building with extensive underground galleries holds more than 40,000 portable artifacts, 2000 of which are on display. Artifacts were moved here from a few other former museums in the region. The range is the entire prehistoric Troad. Displays are multi-lingual. In many cases the original contexts are reproduced.
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+
135
+ Some of the most notable artifacts uncovered at Hisarlik are known as Priam's Treasure. Most of these pieces were crafted from gold and other precious metals. Heinrich Schliemann put this assemblage together from his first excavation site, which he thought to be the remains of Homeric Troy. He gave them this name after King Priam, who is said in the ancient literature to have ruled during the Trojan War. However, the site that housed the treasure was later identified as Troy II, whereas Priam's Troy would most likely have been Troy VIIa (Blegen) or Troy VIi (Korfmann).[note 14] One of the most famous photographs of Sophia made not long after the discovery depicts her wearing a golden headdress, which is known as the "Jewels of Helen" (see under Schliemann above).
136
+
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+ Other pieces that are a part of this collection are:
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+
139
+ Literary Troy was characterized by high walls and towers, summarized by the epithet "lofty Ilium."[80] Some other epithets were "well-walled," "with lofty gates," "with fine towers."[81] Any archaeological candidate for being the literary city would therefore have to show evidence for the walls and towers. Schliemann's Troy fits this qualification very well. High walls and towers are in evidence at every hand. Hisarlik, the name of the hill on which Troy is situated, is Turkish for "the fortress."[82]
140
+
141
+ The walls of Troy, first erected in the Bronze Age between at least 3000 and 2600 BC, were its main defense, as is true of almost any ancient city of urban size. Whether Troy Zero featured walls is not yet known. Some of the known walls were placed on virgin soil (see the archaeology section below). The early date of the walls suggests that defense was important and warfare was a looming possibility right from the beginning.
142
+
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+ The walls surround the citadel, extending for several hundred meters, and at the time they were built were over 17 feet (5.2 m) tall.[83] They were made of limestone, with watchtowers and brick ramparts, or elevated mounds that served as protective barriers.[83]
144
+
145
+ The second run of excavations, under Korfmann, revealed that the walls of the first run were not the entire suite of walls for the city, and only partially represent the citadel. According to Korfmann, "There was also a lower city that went with the Late Bronze Age Troja,...1750-1200 BCE."[84] This city had a perimeter 0f 2.5 kilometres (1.6 mi) end enclosed an area 16 times that of the citadel. It was protected by a ditch surmounted by a wall of mud brick and wood.[85] Moreover, the citadel walls were surmounted by structures of mud brick. The stone part of the walls currently in evidence were "...five meters thick and at least eight meters high - and over that a mudbrick superstructure several meters high...," which totals to about 15 metres (49 ft) for the citadel walls at about the time of the Trojan War. The present-day walls of Troy, then, portray little of the ancient city's appearance, any more than bare foundations characterize a building.
146
+
147
+ Troy I tower and wall
148
+
149
+ South gate wall and tower, Early Troy I through Middle Troy II[86]
150
+
151
+ Troy IV wall
152
+
153
+ Troy VI east tower
154
+
155
+ Troy VI cul-de-sac at east gate
156
+
157
+ Troy VI east tower and wall of cul-de-sac
158
+
159
+ Troy VI wall on the left, Troy IX wall on right. It extends the east gate Troy VI wall on right
160
+
161
+ Troy VI cul-de-sac
162
+
163
+ What Schliemann actually found as he excavated the hill of Hisarlik somewhat haphazardly were contexts in the soil parallel to the contexts of geologic layers in rock. Exposed rock displays layers of a similar composition and fossil content within a layer discontinuous with other layers above and below it. The layer represents an accumulation of detritus over a continuous time, different from the times of the other layers.
164
+
165
+ Similarly Schliemann found layers of distinctive soil each containing more or less distinctive artifacts differing often markedly from other layers. He had no ready explanation for the discontinuity between layers, such as "destruction," although this interpretation has sometimes been applied. Presumably "destruction" is to be interpreted to mean some sort of malicious event perpetrated by humans or a natural disaster, such as an earthquake. In most cases no such disaster can be proved. On the contrary, the "many layers illustrate the gradual development of civilization in northwestern Asia Minor."[83][note 15]
166
+
167
+ The discontinuities of culture in different layers might be explained in a number of ways. A settlement might have been abandoned for peaceful reasons, or it might have undergone a renovation phase. These are hypotheses that must be ruled in or ruled out by evidence, or simply be left unruled until evidence should be discovered.
168
+
169
+ What Schliemann found is that the area now called "the citadel" or "the upper city" was apparently placed on virgin soil. It was protected by fortifications right from the start. The layering effect was caused in part by the placement of new fortifications and new houses over the old. Schliemann called these fortified enclosures "cities" (rightly or wrongly). In his mind the site was composed of successive cities. Like everyone else, he speculated whether a new city represented a different population, and what its relationship to the old was. He numbered the cities I, II, etc., I being on the bottom. Subsequent archaeologists turned the "cities" into layers (rightly or wrongly), named according to the new archaeological naming conventions then being developed.
170
+ The layers of ruins in the citadel at Hisarlik are numbered Troy I – Troy IX, with various subdivisions.
171
+
172
+ Until the late 20th century, these layers represented only the layers on the hill of Hisarlik. Archaeologists following Schliemann picked up the trail of his researches adopting the same fundamental assumptions, culminating in the work and writings of Carl Blegen in the mid-20th century. In a definitive work, Troy and the Trojans, he summarized the layers names and the dates he had adopted for them.[88] Without further excavation, Blegen's was the last word. There were, however, some persistent criticisms not answered to general satisfaction. Hisarlik, about the size of a football field, was not large enough to have been the mighty city of history. It was also far inland, yet the general historical tradition suggested it must have been close to the sea.
173
+
174
+ The issues finally devolved on the necessity for further excavation, which was undertaken by Korfmann starting 1988. He concentrated on the Roman city, which was not suspected as being over Bronze Age remains. A Bronze Age city, at low elevations, was discovered beneath it. As it is unlikely that there were two Troys side by side, the lower city must have been the main seat of residence, to which the upper city served as citadel. Korfman now referred to the layers of the lower city as associated with the layers of the citadel. The same layering scheme was applicable. The lower city was many times the size of the citadel, answering the size objection.
175
+
176
+ Meanwhile independent geoarchaeological research conducted by taking ground cores over a wide area of the Troad were demonstrating that, in the time of Troy I, "... the sea was right at the foot of 'Schliemann's Trench' during the earliest periods of Troja."[89] A few thousand years earlier the ridge of Troy was partly surrounded by an inlet of the sea occupying the now agricultural area of the lower Scamander River. Troy was founded as an apparently maritime city on the shore of this inlet, which persisted throughout the early layers and was present to a lesser degree, farther away, subsequently. The harbor at Troy, however, was always small, shallow, and partially blocked by wetlands. It was never a "great harbor" able to collect maritime traffic through the Dardanelles.[90] The current water table depends on the degree of irrigation of the now agricultural lands. Trench flooding has slowed investigation of the lower levels in the lower city.
177
+
178
+ The whole course of archaeological investigation at Troy has resulted in no single chronological table of layers. Moreover, due to limitations on the accuracy of C14 dating, the tables remain relative; i.e., absolute, or calendar dates, cannot be precisely assigned. In regions of the Earth where both history and C14 dating are available, there is often a gap between them, termed by Renfrew a chronological or archaeological "fault line." The two models, historical and archaeological, do not correspond, just as the contexts on either side of a geologic fault line do not correspond. "This line divides all Europe except the Aegean from the Near East."[91]
179
+
180
+ The table below concentrates on two systems of dates: Blegen's from Troy and the Trojans,[88][note 16], representing the last of the trend from Schliemann to the mid-20th century, and Korfmann's, from Troia in Light of New Research in the early years of the 21st century, after he had had a chance to establish a new trend and new excavations.[92]
181
+
182
+ Prior to Korfmann's excavations, the nine-layer model was considered comprehensive of all the material at Troy. Korfmann discovered that the city was not placed on virgin soil, as Schliemann had concluded. There is no reason not to think that, in the areas he tested, Schliemann did find that Troy I was on virgin soil. Korfmann discovered a layer previous to Troy I under a gate to Troy II. He dated it 3500 BC to 2920 BC, but did not assign a name. The current director of excavation at Troy, Rüstem Aslan, is calling it Troy 0 (zero).[93] Roman numerals have no zero, but zero is one number less than I.
183
+
184
+ Troy 0 has been omitted from the table below, due to the uncertainty of its general status. Archaeologists at the site before Korfmann had thought that Troy I began with the Bronze Age at 3000 BC. Troy zero is before this date. The remains of the layer are not very substantial. Whether the layer is to be counted as part of the preceding Chalcolithic, or whether the dates of the Bronze Age are to be changed, has not been decided through the regular channel of journal articles. One 2016 PhD Thesis complained: "... the stratigraphic sequence of the renewed excavations is presented differently by different collaborators of Korfmann ... So, until an agreed stratigraphy of Korfmann’s cycle is published, the employment of Troy as a yardstick for the whole of the Anatolian EBA remains problematic."[94][note 17]
185
+
186
+ For example, in Korfmann 2003, p. 31 Korfmann elaborates beyond the chronology of Cobet's table to make new proposals regarding the layer, Troy VIIa (which he also presents in the Guidebook): "Troia VIIa should be assigned culturally to
187
+ Troia VI," asserting that "there were no substantial differences in the material culture between the two periods." He suggests that Dőrpfeld's classification of the material subsequently in VIIa as VIi should be restored, claiming that, regarding the details, Blegen had been "entirely in agreement" even though his chronology featured Troy VIIa.[note 14] He then laments "the old terminology has, unfortunately, been retained. Confusion is to be avoided at all costs." As this view has not yet been tested in the journals and is not universal, it is mainly omitted from the table (Cobet's chart, however, includes Korfmann's VIIb 3.) This new and yet unresolved material, including Troy Zero, may, however, be included in the sections and links below reporting on specific layers
188
+
189
+ Korfmann also found that Troy IX was not the end of the settlements. Regardless of whether the city was abandoned at 450 AD, a population was back for the Middle Ages, which, for those times, was under the Byzantine Empire. As with Troy Zero, no conventional scholarly classification has been tested in the journals. The literature mentions Troy X, and even Troy XI, without solid definition. The table below therefore omits them.
190
+
191
+ The sequence of archaeological layering at one site evidences the relative positions of the corresponding periods at that site; however, these layers often have a position relative to periods at other sites. It is possible to define relative periods over a wide region of sites and for a larger slice of time. Determining wider correspondences is a major objective of archaeology. The establishment of a "yardstick," or reliable sequence, such as the elusive one mentioned above, is a desirable outcome of archaeological analysis.
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+
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+ The table below states the broader connections under "General Period." It references primarily the chronologies presented in the educational site created and maintained by Jeremy Rutter and team and published by Dartmouth College, entitled Aegean Prehistoric Archaeology.[95][note 18] The time period is generally "the Bronze Age," which has an early (EB or EBA), a middle (MB or MBA), and a late (LB or LBA). The sites are distributed over Crete ("Minoan," or M), the Cyclades ("Cycladic," or C), the Greek mainland ("Helladic," or H), and Western Turkey ("Western Anatolian," no abbreviation).
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+
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+ The first city on the site was founded in the 3rd millennium BC. During the Bronze Age, the site seems to have been a flourishing mercantile city, since its location allowed for complete control of the Dardanelles, through which every merchant ship from the Aegean Sea heading for the Black Sea had to pass. Cities to the east of Troy were destroyed, and although Troy was not burned, the next period shows a change of culture indicating a new people had taken over Troy.[96] The first phase of the city is characterized by a smaller citadel, around 300 ft in diameter, with 20 rectangular houses surrounded by massive walls, towers, and gateways.[83] Troy II doubled in size and had a lower town and the upper citadel, with the walls protecting the upper acropolis which housed the megaron-style palace for the king.[97] The second phase was destroyed by a large fire, but the Trojans rebuilt, creating a fortified citadel larger than Troy II, but which had smaller and more condensed houses, suggesting an economic decline.[83] This trend of making a larger circuit, or extent of the walls, continued with each rebuild, for Troy III, IV, and V. Therefore, even in the face of economic troubles, the walls remained as elaborate as before, indicating their focus on defense and protection.
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+
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+ When Schliemann came across Troy II, in 1871, he believed he had found Homer's city. Schliemann and his team unearthed a large feature he dubbed the Scaean Gate, a western gate unlike the three previously found leading to the Pergamos.[98] This gate, as he describes, was the gate that Homer had featured. As Schliemann states in his publication Troja:
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+ "I have proved that in a remote antiquity there was in the plain of Troy a large city, destroyed of old by a fearful catastrophe, which had on the hill of Hisarlık only its Acropolis with its temples and a few other large edifices, southerly, and westerly direction on the site of the later Ilium; and that, consequently, this city answers perfectly to the Homeric description of the sacred site of Ilios."[99]Also, he uncovered what he referred to as The Palace of Priam, after the king during the Trojan War.[100] This reference is incorrect because Priam lived nearly a thousand years after Troy II.
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+
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+ Troy VI was destroyed around 1250 BC, probably by an earthquake. Only a single arrowhead was found in this layer, and no remains of bodies. However, the town quickly recovered and was rebuilt in a layout that was more orderly. This rebuild continued the trend of having a heavily fortified citadel to preserve the outer rim of the city in the face of earthquakes and sieges of the central city.[97]
201
+
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+ Troy VI can be characterized by the construction of the pillars at the south gate. There appears to be no structural use for the pillars. The pillars have an altar-like base and an impressive magnitude. This provides some clues, and they most likely were used as a symbol for the religious cults of the city.[101] Another characterizing feature of Troy VI is the tightly packed housing near the Citadel and construction of many cobble streets. Although only few homes could be uncovered, this is due to reconstruction of Troy VIIa over the tops of them.[102]
203
+
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+ Also, discovered in 1890, in this layer of Troy VI was Mycenaean pottery. This pottery suggests that during Troy IV, the Trojans still had trade with the Greeks and the Aegean. Furthermore, there were cremation burials discovered 400m south of the citadel wall. This provided evidence of a small lower city south of the Hellenistic city walls. Although the size of this city is unknown due to erosion and regular building activities, there is significant evidence that was uncovered by Blegen in 1953 during an excavation of the site. This evidence included settlements just above bedrock and a ditch thought to be used for defense. Furthermore, the small settlement itself, south of the wall, could have also been used as an obstacle to defend the main city walls and the citadel.[103]
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+
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+ The topic still under debate is whether Troy was primarily an Anatolian-oriented or Aegean-oriented metropolis. While it is true that the city would have had a presence in the Aegean, pottery finds and architecture strongly hint at an Anatolian orientation. Only about one percent of the pottery discovered during excavation of Troy VI was Mycenaean. The large walls and gates of the city are closely related to many other Anatolian designs. Furthermore, the practice of cremation is Anatolian. Cremation is never seen in the Mycenaean world. Anatolian hieroglyphic writing along with bronze seals marked with Anatolian hieroglyphic Luwian were also uncovered in 1995. These seals have been seen in approximately 20 other Anatolian and Syrian cities from the time (1280 - 1175 BC).[104]
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+ Still, Troy VI was dominated by long distance trade. Troy VI during the height of its establishment held anywhere from 5,000 to 10,000 people. At its time, Troy would have been a large and significant city.[104] The location of Troy was extremely practical in the Early Bronze Age (2000–1500 BC). It acted as a middle ground for long distance trade with regions as far distant as Afghanistan, the Persian Gulf, the Baltic region, Egypt, and the western Mediterranean in the Middle and Late Bronze Ages. Earlier trade connections during the Early and Middle Bronze Ages provided Troy VI with favorable power in the long distance trade industry of the region. The amount of objects thought to be going through Troy VI would have been quite large, obtaining metals from the east and various objects from the west including perfumes and oils. This is known due to the findings of hundreds of shipwrecks off the Turkish coast. Found in these ships was an abundance of goods. Some of these ships carried over 15 tons in goods. The goods discovered in these wrecks included copper ingots, tin ingots, glass ingots, bronze tools and weapons, ebony and ivory, ostrich egg shells, jewelry and large amounts of pottery from across the Mediterranean.[105]
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+ There have been 210 shipwrecks discovered in the Mediterranean from the Bronze Age. Of these 210, 63 were discovered off the Turkish coastline. This provides a great deal of evidence for Troy VI being a prominent trading center for the region. But, the evidence at the site of Troy itself is minimal. Looking at the layers of Troy VI, we discover that there is little documentation of the excavation of this layer, and little documentation of the goods discovered in this layer. We also know that there were few trading centers during the Late Bronze Age. This is due to the low volume of trade during this period. The trading centers would have most likely been directly along trade routes. Troy is just north of most major long-distance trade routes. It may be unfair to classify Troy VI as a trading center but we do know that Troy VI was a prominent metropolis that did contribute to the trade of the region.[105]
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+ Troy VIIa can be highlighted by most of the population of Troy moving within the walls of the Citadel. This is most likely due to the threat from the Mycenaeans.[106] Troy VI is believed to have been destroyed by an earthquake. This would not have been uncommon. Earthquakes are common throughout the region. Troy VIIa is believed to be built over the ruined Troy VI, which makes the excavation process of Troy difficult.[107]
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+ Troy VIIa, which has been dated to the mid-to-late-13th century BC, is the most often cited candidate for the Troy of Homer. Troy VIIa appears to have been destroyed by war.[108] The evidence of fire and slaughter around 1184 BC, which brought Troy VIIa to a close, led to this phase being identified with the city besieged by the Greeks during the Trojan War. This was immortalized in the Iliad written by Homer.[109]
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+ Initially, the layers of Troy VI and VII were overlooked entirely, because Schliemann favoured the burnt city of Troy II. It was not until the need to close "Calvert's Thousand Year Gap" arose—from Dörpfeld's discovery of Troy VI—that archaeology turned away from Schliemann's Troy and began working towards finding Homeric Troy once more.[110]
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+ "Calvert's Thousand Year Gap" (1800–800 BC) was a period not accounted for by Schliemann's archaeology and thus constituted a hole in the Trojan timeline. In Homer's description of the city, a section of one side of the wall is said to be weaker than the rest.[111] During his excavation of more than three hundred yards of the wall, Dörpfeld came across a section very closely resembling the Homeric description of the weaker section.[112] Dörpfeld was convinced he had found the walls of Homer's city, and now he would excavate the city itself. Within the walls of this stratum (Troy VI), much Mycenaean pottery dating from Late Helladic (LH) periods III A and III B (c. 1400–c. 1200 BC) was uncovered, suggesting a relation between the Trojans and Mycenaeans. The great tower along the walls seemed likely to be the "Great Tower of Ilios".[113]
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+ The evidence seemed to indicate that Dörpfeld had stumbled upon Ilios, the city of Homer's epics. Schliemann himself had conceded that Troy VI was more likely to be the Homeric city, but he never published anything stating so.[114] The only counter-argument, confirmed initially by Dörpfeld (who was as passionate as Schliemann about finding Troy), was that the city appeared to have been destroyed by an earthquake, not by men.[115] There was little doubt that this was the Troy of which the Mycenaeans would have known.[116]
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+ The archaeologists of Troy concerned themselves mainly with prehistory; however, not all the archaeology performed there falls into the category of prehistoric archaeology. Troy VIII and Troy IX are dated to historical periods. Historical archaeology illuminates history. In the LBA records mentioning Troy begin to appear in other cultures. This type of evidence is termed protohistory. The literary characters and events must be classified as legendary. Prehistoric Troy is also legendary Troy. The legends are not history or protohistory, as they are not records. It was the question of their historicity that attracted the interest of such archaeologists as Calvert and Schliemann. After many decades of archaeology, there are still no answers. There is still a "fault line" between history or legend and archaeology.
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+ If Homeric Troy is not a fantasy woven in the 8th century by Greek oral poets passing on a tradition of innovating new poems at festivals, as most archaeologists hoped it was not, then the question must be asked, what archaeological level represents Homeric Troy? Only two credible answers are available, which are the same answer: Troy VIIa in the Blegen scheme,[117] identical to Troy VIi in the scheme suggested by Korfmann.[118] After an earthquake brought down the walls of the city at its floruit about 1300 BC, the same people rebuilt the city even more magnificent than before. This event is considered the start of the LBA, and Homeric Troy is considered to be LBA Troy.
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+ Both Blegen and Korfmann endorse a starting date of about 1300 BC. Blegen has it ending early at 1260 BC,[note 19] but Korfmann runs it up to 1190 BC (or 1180 BC elsewhere). He abolishes VIIa, and substitutes for it VIi, more in keeping with the splendor of VI; after all, they were the same people. He estimates the population at 10,000.[119] The end of the period is marked by weapons left laying around, skeletons, and burnt objects, considered the result of the Trojan War. Coincidentally this is the very period referenced by Egyptian and Hittite records of Troy. They hold out some hope of a protohistorical connection.
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+ In the 1920s, the Swiss scholar Emil Forrer proposed that the placenames Wilusa and Taruisa found in Hittite texts should be identified with Ilion and Troia, respectively.[120] He further noted that the name of Alaksandu, a king of Wilusa mentioned in a Hittite treaty, is quite similar to Homer's Paris, whose birthname was Alexandros. These identifications were rejected by many scholars as being improbable or at least unprovable. However, Trevor Bryce championed them in his 1998 book The Kingdom of the Hittites, citing a piece of the Manapa-Tarhunda letter referring to the kingdom of Wilusa as beyond the land of the Seha River (the classical Caicus and modern Bakırçay) and near the land of "Lazpa" (Lesbos).
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+ The excavation of the lower city uncovered a water distribution system containing 160 metres (520 ft) of tunnels tapping sources higher up on the ridge. Dates from the floor deposits obtained by the Uranium-thorium dating method indicate that water was flowing through the tunnels "as early as the third millenium BC;" thus the early city made sure that it had an internal water supply.[121] In 1280 BC a treaty between the reigning monarchs of the Hittite and Trojan states, Muwatalli II and Alaksandu of Wilusa respectively, invoked the water god, KASKAL_KUR, who was associated with an underground tunnel, adding weight to the theory that Wilusa is identical to archaeological Troy.
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+ Among the documents mentioning Troy are the Tawagalawa letter (CTH 181) was found to document an unnamed Hittite king's correspondence to the king of the Ahhiyawa, referring to an earlier "Wilusa episode" involving hostility on the part of the Ahhiyawa. The Hittite king was long held to be Mursili II (c. 1321–1296), but, since the 1980s, his son Hattusili III (1265–1240) is commonly preferred, although his other son Muwatalli (c. 1296–1272) remains a possibility.
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+ Inscriptions of the New Kingdom of Egypt also record a nation T-R-S as one of the Sea Peoples who attacked Egypt during the XIX and XX Dynasties. An inscription at Deir el-Medina records a victory of Ramesses III over the Sea Peoples, including one named "Tursha" (Egyptian: [twrš3]). It is probably the same as the earlier "Teresh" (Egyptian: [trš.w]) on the stele commemorating Merneptah's victory in a Libyan campaign around 1220 BC.
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+ The identifications of Wilusa with Troy and of the Ahhiyawa with Homer's Achaeans remain somewhat controversial but gained enough popularity during the 1990s to be considered majority opinion. That agrees with metrical evidence in the Iliad that the name ᾽Ιλιον (Ilion) for Troy was formerly Ϝιλιον (Wilion) with a digamma.[note 3]
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+ From the beginning of the archaeology, the question of what language was spoken by the Trojans was prominent. Various proposals were made, but they remained pure speculation. No evidence seemed to have survived whatever. That they might be Greek was considered. However, if they were, the question of why they were not in the Achaean domain, but were opposed to the Achaeans, was an even greater mystery. Passages from the Iliad suggested that, not only were the Trojans not Greek, but the army defending Troy was composed of different language speakers arrayed by nationality.[122]
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+ Finally in the middle of the 20th century Linear B was deciphered and a large number of documentary tablets were able to be read. The language is an early dialect of Greek, even earlier than the Homeric dialect. Many Greek words were in the early stage of formation. The digamma abounds. Linear B tablets have been found at the major centers of the Achaean domain. None, however, come from Troy.
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+ The documents in Linear B basically inventory the assets of Mycenaean palace-states: foods, textiles, ceramics, weapons, lands, and above all manpower, especially people held in some sort of servitude. Civilizations of the times were slave societies. The terms of servitude, however, varied widely. A study by Efkleidou in 2004 detailed the types of servitude mentioned in the Linear B tablets. To her way of thinking, the main elements of servitude are that servants are outsiders, not part of the customary social structure, and that they are coerced into their positions. Someone has authority over them, whom she calls a "superior," designated in Greek by the genitive case: "servant of ...." One of the categories of Mycenaean servant is the do-er-o (masc) and do-er-a (fem), Greek doulos, pl douloi, and doula, pl doulai.[123] A specific type of doulos is the te-o-jo do-er-o, theoio doulos, "servant of God," a temple assistant of some sort, whose superior was the deity. These two categories were not badly off, being palace artisans, and receiving land for their services. In addition were the ra-wi-ja-ja, the lawiaiai, "captives." These were kept in groups and performed what would be termed today "factory work." The tablets, being ephemeral in nature, do not always classify these types, but they are detectable from the naming conventions, or lack of them, and the type of work. Efkleidou uses the term "dependent." In all she tallied 5233 dependents in the tablets.
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+ Perhaps most relevant to the times are named groups of women, the group name being an ethnic or a craft name. One such group called just "captives" gives a hint to their class of servitude. The ethnic names show that western Anatolia and the islands off it are being favored.[note 20] Other groups implicitly from the region are named after the type of work they do, especially the textile workers: finishers, spinners, and a group of flax handlers (ri-ne-ja, or lineiai) composed of 82 women with 61 female children and 56 male. Other groups are male bronzesmiths, house and ship builders.[124] The majority of the females were textile workers, a development foreshadowed in the initial scene of the Iliad, in which the priest, Chryses, entreats Agamemnon to ransom his daughter, Chryseis, only to be refused with the statement that she would be frequenting his bed and working his loom far away in Argos.
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+ In the tablets, the coast of Anatolia is under attack by Mycenaean centers of the Achaeans, especially the center at Pylos (pu-ro). Since the tablets, which were manufactured ad hoc of fresh clay and immediately engraved with writing, only survived by being baked in the fires that destroyed the palaces, their dates depend on the those dates of destruction. The Pylos tablets record the dispatch of a fleet of "rowers" and soldiers under a "commander" to the Gulf of Corinth, and then the palace is gone, burned in its own oil. If pu-ro is the Homeric Pylos, then the date is after the Trojan War, as the legendary Pylos survived it intact. None of the names of the important men at these centers are anything like the names of the Homeric legends. Presumably, the latter had all died in their time and had been replaced by men unknown to legend, but profiting from the fall of Troy. A second possibility would be that the legends are totally imaginary, contrary to the hopes and expectations of the first archaeologists.
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+ This time between the Trojan War and the burning of the palaces fits into another historical period, the time of the Sea Peoples. These were ethnicities from Achaea, Dardania, Etruria, Sicilia, Sardinia, and elsewhere, who, abandoning the norms of civilization, took to a life of marauding and piracy, disrupting trade, transportation, peace, and security. They placed colonies as bases.[note 21] The eastern Mediterranean became a wilderness. Cities withdrew from the coast. Isolation set in.
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+ Surprisingly, Trojan names began turning up in the archives of ancient Pylos, capital of the legendary Nestor. They were of persons kept in a servile capacity, from which the universal conclusion was that they were descended from slaves taken at Troy. Etymological analysis by linguists revealed that they were not native Greek names, suggesting that the Trojans were not Greek.
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+ A theory began to gain influence based on the Aeneid that the Trojans were Etruscan.[note 22] This theory purported to explain how the Etruscans arrived in Italy. During the 20th century, however, Etruscan archaeology investigated thousands of Etruscan sites over most of Italy, except for the Greek regions in south Italy and the Italic regions of central Italy. Moreover, Etruscan inscriptions were found in at least one valley leading to a pass over the Alps. The sites dated as early as the Bronze Age. It was soon clear that the theory of a general Etruscan migration from Troy to most of Italy was the least likely scenario.[note 23] The Etruscan theory died slowly. Its advocates looked for hidden pockets of Etruscans in the backlands of Anatolia and looked for hope in some shallow genetic studies purporting to relate the inhabitants of Tuscany to the inhabitants of Turkey.[note 24]
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+ Meanwhile a greater question came to the fore. Throughout the Bronze Age the greatest power in Anatolia was the Hittites, with capital in central Anatolia. Why were there no links to them? How could the coastal states have avoided them? Anatolian studies expanded in the late 20th century. Those states had not avoided them, they were subject to them. Previously unknown scripts were found to be in Anatolian languages. The dominant one on the coast was Luwian. In the Luwian range west of the Hittite capital there was no room for any Etruscans. Whatever he was, Aeneas was not Etruscan, and whatever the ancestry of the imperial family at Rome was, which knew Etruscan and was counted as Tuscan, it derived no authority from ancient Troy.
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+ The 1995 discovery of a Luwian biconvex seal at Troy sparked heated debate over the language that was spoken in Homeric Troy. Frank Starke of the University of Tübingen argued that the name of Priam, king of Troy at the time of the Trojan War, is related to the Luwian compound Priimuua, which means "exceptionally courageous".[125] Starke adds: "The certainty is growing that Wilusa/Troy belonged to the greater Luwian-speaking community," although it is not entirely clear whether Luwian was primarily the official language or in daily colloquial use.[126]. The tablet was discovered in the lower city, archaeologically out of the way until now, but undoubtedly more populous and frequented than the citadel. It is possible that the major archive site has yet to be discovered at Troy, if any survived.
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+ B.W. Fortson, IV, defines the Greek Dark Ages as "The period from the demise of Mycenaean civilization to the earliest appearance of alphabetic Greek in the eighth century ...."[127] The idea is that in the breakdown of peace and stability suffered by the Mycenaean kingdoms they entered a period of fear, isolation, and economic depression in which their writing system was lost causing a subsequent deficit of written records, interpreted by later historians as "darkness." The view is too simple, however.
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+ While it is true that the palaces were destroyed by fire, it is untrue that they were all burned in the same year or even the same decade by a single wave of Dorian tribes from the region later known as Macedonia. The dates of the destructions differ by as much as a generation. Chadwick asks, "... where were all the Dorians during the Mycenaean period? And why were they content to wait in the wings until the time was right for this intrusion?"[128] His own theory was that the Mycenaeans were incendiary to each other's palaces in a rash of infradynastic conflicts. These would have occupied the entire 11th century BC. There was no sudden influx of all the Dorians in one great invasion, but rather an insistent occupation of the Peloponnesus over a century or more. It has to be counted as Dorian from the 10th century BC on. Most of the former Achaean inhabitants escaped to the now depopulated coast of Anatolia as Ionians and Aeolians. Athens remained firm.
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+ Among the Achaeans of Cyprus, on the edge of the Greek world, writing failed to disappear. They continued to write their own conservative dialect, Arcadocypriot Greek, in a few scripts of Cypriote syllabary, which they had innovated on the model of Linear A and Linear B. They were fairly isolated from their former homeland by the spread of Dorians to Crete, the southern Cyclades, and southern Anatolia. When the concept of a Greek alphabet arrived, they innovated with the Phoenician alphabet to make it fit their language, and the two systems continued side-by-side until Hellenistic times, when Attic became the common dialect. Meanwhile their dialect continued in the hills of Arcadia, but it had no writing system there. This dark age interlude in Greece is not generally interpreted as a return to prehistoric times. It is a historic age with gaps in its history, which is how the archaeologists treat it.
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+ In both Blegen and Korfmann, Homeric Troy ends at the end of Blegen's Troy VIIa, Korfmann's VIi, with the burning of the city, presumably by Achaeans. Legend has the Trojans vanishing away, either escaping, as did Aeneas and his very large band, being slaughtered, as were Priam and his wife, or being carted off into slavery, as were the literary Trojan women. Apparently, no Trojans seem to have been left. Their enemies would have cleared them entirely away, leaving the ruined city vacant and non-dangerous.
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+ The archaeology suggests that the literary implication of a deserted city is probably not true. After a suitable interval of hiding somewhere else in the region, perhaps with the Dardanians, who were not defeated, but appeared as marauders among the sea peoples, or further inland with the Hittites, the Trojan remnants returned to Troy to rebuild Troy VIIb, which, according to Blegen, "... obviously represents a direct survival of the culture that prevailed in Troy VIIa."[129] The initial VIIb period is VIIb1, which Korfmann suggests should be VIj.[130] and regards as "transitional to the Early Iron Age." As yet, however, it is contemporary with LHIIIC (LBA) pottery on the mainland. The reconstruction does not appear to have been opposed by the palaces, such as at Pylos, which were still standing. The return to a simpler pottery causes Korfmann to hypothesize a "humble folk" investment of the ruins.
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+ Troy VIIb2 begins contemporaneously with LHIIIC, but at about 1050 BC the last of IIIC disappears, to be replaced by Sub-Mycenaean pottery, a short-lived Mycenaean-like pottery with geometric motifs, considered transitional to Geometric pottery, the ware characteristic everywhere in the Greek world of the Dark Age. The palaces can be counted as vanished, as the last pottery at Pylos was LH IIIC. Apparently, the city of the "humble Trojans" could not maintain itself, but was overrun or replaced. The latter part of Troy IIIb2 sees the replacement of their pottery with wares, such as "Knobbed Ware," characteristic of the Balkan-Black Sea region. The Luwian seal presents a problem, as it is dated Troy VIIb2. Luwian speakers would not have been as far away as the northern Black Sea. If the seal is from early VIIb2, however, it can represent the last of the Luwian speakers at Troy. A mixed culture was certainly possible. Priam's wife, Hecuba, had been a Phrygian.
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+ Blegen ends his tale of Troy VII with VIIb2 around 1100 BC. The city was burned one last time, an event contemporaneous with the general destruction of the Mycenaean palaces. This would be the ethnical end of the Trojans at Troy by abandonment, but Blegen has a final suggestion. Troy VI was characterized by what Blegen calls "Grey Minyan Ware," now Anatolian Minyan ware.[131][note 25] or Anatolian Grey Ware. After the abandonment of the city, the ware appears in the highlands, leading Blegen to conjecture that the Trojans gradually withdrew in that direction.
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+ The more recent excavations turned up additional information. In the lower city was pottery from the early and middle Proto-geometric period, characteristic of the Dark Age. The Trojans may have escaped to the hills, but their burned city was occupied by their incendiary opponents, whoever they were. They brought iron with them, relying on the superior strength of iron weapons for their victory. Korfmann creates a new period for them, Troy VIIb3, 1020-950 BC.
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+ For reasons unknown, the Iron-age people left their settlement about 950 BC, leaving it abandoned. Korfmann calls this interval a hiatus, meaning of residential occupation. A Greek colony arrived there to plant a new city about 750 BC, archaeological Troy VIII. They leveled the top of the mound to construct a temple to Athena, thus identifying themselves as being in the Attic-Ionic culture, as opposed to the Aeolic Greeks (Boeotia) who had previously been settling the north coast of Anatolia. The leveling process destroyed the previous structures at the center of the citadel. As Homeric Troy had been called "sacred Ilium," Korfmann asserts that a temple district may have been maintained there during the apparent abandonment period, but whose is not known.
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+ In 480 BC, the Persian king Xerxes sacrificed 1,000 cattle at the sanctuary of Athena Ilias while marching through the Hellespontine region towards Greece.[132] Following the Persian defeat in 480–479, Ilion and its territory became part of the continental possessions of Mytilene and remained under Mytilenaean control until the unsuccessful Mytilenean revolt in 428–427. Athens liberated the so-called Actaean cities including Ilion and enrolled these communities in the Delian League. Athenian influence in the Hellespont waned following the oligarchic coup of 411, and in that year the Spartan general Mindaros emulated Xerxes by likewise sacrificing to Athena Ilias. From c. 410–399, Ilion was within the sphere of influence of the local dynasts at Lampsacus (Zenis, his wife Mania, and the usurper Meidias) who administered the region on behalf of the Persian satrap Pharnabazus.
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+ In 399, the Spartan general Dercylidas expelled the Greek garrison at Ilion who were controlling the city on behalf of the Lampsacene dynasts during a campaign which rolled back Persian influence throughout the Troad. Ilion remained outside the control of the Persian satrapal administration at Dascylium until the Peace of Antalcidas in 387–386. In this period of renewed Persian control c. 387–367, a statue of Ariobarzanes, the satrap of Hellespontine Phrygia, was erected in front of the temple of Athena Ilias.[133] In 360–359 the city was briefly controlled by Charidemus of Oreus, a Euboean mercenary leader who occasionally worked for the Athenians.[134] In 359, he was expelled by the Athenian Menelaos son of Arrabaios, whom the Ilians honoured with a grant of proxeny—this is recorded in the earliest civic decree to survive from Ilion.[135] In May 334 Alexander the Great crossed the Hellespont and came to the city, where he visited the temple of Athena Ilias, made sacrifices at the tombs of the Homeric heroes, and made the city free and exempt from taxes.[136] According to the so-called 'Last Plans' of Alexander which became known after his death in June 323, he had planned to rebuild the temple of Athena Ilias on a scale that would have surpassed every other temple in the known world.[137]
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+
280
+ Antigonus Monophthalmus took control of the Troad in 311 and created the new city of Antigoneia Troas which was a synoikism of the cities of Skepsis, Kebren, Neandreia, Hamaxitos, Larisa, and Kolonai. In c. 311–306 the koinon of Athena Ilias was founded from the remaining cities in the Troad and along the Asian coast of the Dardanelles and soon after succeeded in securing a guarantee from Antigonus that he would respect their autonomy and freedom (he had not respected the autonomy of the cities which were synoikized to create Antigoneia).[138] The koinon continued to function until at least the 1st century AD and primarily consisted of cities from the Troad, although for a time in the second half of the 3rd century it also included Myrlea and Chalcedon from the eastern Propontis.[139] The governing body of the koinon was the synedrion on which each city was represented by two delegates. The day-to-day running of the synedrion, especially in relation to its finances, was left to a college of five agonothetai, on which no city ever had more than one representative. This system of equal (rather than proportional) representation ensured that no one city could politically dominate the koinon.[140] The primary purpose of the koinon was to organize the annual Panathenaia festival which was held at the sanctuary of Athena Ilias. The festival brought huge numbers of pilgrims to Ilion for the duration of the festival as well as creating an enormous market (the panegyris) which attracted traders from across the region.[141] In addition, the koinon financed new building projects at Ilion, for example a new theatre c. 306 and the expansion of the sanctuary and temple of Athena Ilias in the 3rd century, in order to make the city a suitable venue for such a large festival.[142]
281
+
282
+ In the period 302–281, Ilion and the Troad were part of the kingdom of Lysimachus, who during this time helped Ilion synoikize several nearby communities, thus expanding the city's population and territory.[note 26] Lysimachus was defeated at the Battle of Corupedium in February 281 by Seleucus I Nikator, thus handing the Seleucid kingdom control of Asia Minor, and in August or September 281 when Seleucus passed through the Troad on his way to Lysimachia in the nearby Thracian Chersonese Ilion passed a decree in honour of him, indicating the city's new loyalties.[143] In September Seleucus was assassinated at Lysimachia by Ptolemy Keraunos, making his successor, Antiochus I Soter, the new king. In 280 or soon after Ilion passed a long decree lavishly honouring Antiochus in order to cement their relationship with him.[note 27] During this period Ilion still lacked proper city walls except for the crumbling Troy VI fortifications around the citadel, and in 278 during the Gallic invasion the city was easily sacked.[144] Ilion enjoyed a close relationship with Antiochus for the rest of his reign: for example, in 274 Antiochus granted land to his friend Aristodikides of Assos which for tax purposes was to be attached to the territory of Ilion, and c. 275–269 Ilion passed a decree in honour of Metrodoros of Amphipolis who had successfully treated the king for a wound he received in battle.[145]
283
+
284
+ A new city called Ilium (from Greek Ilion) was founded on the site in the reign of the Roman Emperor Augustus. It flourished until the establishment of Constantinople, became a bishopric in the Roman province Hellespontus (civil Diocese of Asia), but declined gradually in the Byzantine era.
285
+
286
+ The city was destroyed by Sulla's rival, the Roman general Fimbria, in 85 BC following an eleven-day siege.[146] Later that year when Sulla had defeated Fimbria he bestowed benefactions on Ilion for its loyalty which helped with the city's rebuilding. Ilion reciprocated this act of generosity by instituting a new civic calendar which took 85 BC as its first year.[147] However, the city remained in financial distress for several decades, despite its favoured status with Rome. In the 80s BC, Roman publicani illegally levied taxes on the sacred estates of Athena Ilias and the city was required to call on L. Julius Caesar for restitution; while in 80 BC, the city suffered an attack by pirates.[148] In 77 BC the costs of running the annual festival of the koinon of Athena Ilias became too pressing for both Ilion and the other members of the koinon and L. Julius Caesar was once again required to arbitrate, this time reforming the festival so that it would be less of a financial burden.[149] In 74 BC the Ilians once again demonstrated their loyalty to Rome by siding with the Roman general Lucullus against Mithridates VI.[150] Following the final defeat of Mithridates in 63–62, Pompey rewarded the city's loyalty by becoming the benefactor of Ilion and patron of Athena Ilias.[151] In 48 BC, Julius Caesar likewise bestowed benefactions on the city, recalling the city's loyalty during the Mithridatic Wars, the city's connection with his cousin L. Julius Caesar, and the family's claim that they were ultimately descended from Venus through the Trojan prince Aeneas and therefore shared kinship with the Ilians.[152]
287
+
288
+ In 20 BC, the Emperor Augustus visited Ilion and stayed in the house of a leading citizen, Melanippides son of Euthydikos.[153] As a result of his visit, he also financed the restoration and rebuilding of the sanctuary of Athena Ilias, the bouleuterion (council house) and the theatre. Soon after work on the theatre was completed in 12–11 BC, Melanippides dedicated a statue Augustus in the theatre to record this benefaction.[154]
289
+
290
+ No later than the 4th century, it was a suffragan of the provincial capital's Metropolitan Archdiocese of Cyzicus, in the sway of the Patriarchate of Constantinople.
291
+ Several bishops are historically documented:
292
+
293
+ The diocese was nominally restored no later than 1926 as Latin Titular bishopric of Ilium (Latin) / Ilio (Curiate Italian) / Ilien(sis) (Latin adjective).
294
+
295
+ It has been vacant for decades, having had the following incumbents, so far of the fitting Episcopal (lowest) rank:
296
+
297
+ A small minority of contemporary writers argue that Homeric Troy was not at the Hisarlik site, but elsewhere in Anatolia or outside it—e.g. in England,[155] Pergamum,[156] Scandinavia,[157] or Herzegovina.[158] These proposals have not been accepted by mainstream scholarship.
298
+
299
+ Such was the fame of the Epic Cycle in Roman and Medieval times that it was built upon to provide a starting point for various founding myths of national origins. The most influential, Virgil's Aeneid, traces the journeys of the Trojan prince Aeneas, supposed ancestor of the founders of Rome and the Julio-Claudian dynasty. In a later era, the heroes of Troy, both those noted in Homer and those invented for the purpose, often continued to appear in the origin stories of the nations of Early Medieval Europe.[159][160] The Roman de Troie was common cultural ground for European dynasties,[161] as a Trojan pedigree was both gloriously ancient and established an equality with the ruling class of Rome. A Trojan pedigree could justify the occupation of parts of Rome's former territories.[159]
300
+
301
+ Dionysius of Halicarnassus writes that the Trojans were Greek people who were originally from the Peloponnese.[162]
302
+
303
+ On that basis, the Franks filled the lacunae of their legendary origins with Trojan and pseudo-Trojan names: in Fredegar's 7th-century chronicle of Frankish history, Priam appears as the first king of the Franks.[163][full citation needed] The Trojan origin of France was such an established article of faith that in 1714, the learned Nicolas Fréret was Bastilled for showing through historical criticism that the Franks had been Germanic, a sore point counter to Valois and Bourbon propaganda.[164][full citation needed]
304
+
305
+ In similar manner, Geoffrey of Monmouth reworked earlier material such as the Historia Brittonum to trace the legendary kings of the Britons from a supposed descendant of Aeneas called Brutus.
306
+
307
+ Likewise, Snorri Sturluson, in the prologue to his Icelandic Prose Edda, traced the genealogy of the ancestral figures in Norse mythology to characters appearing at Troy in Homer's epic, notably making Thor to be the son of Memnon. Sturluson referred to these figures as having made a journey across Europe towards Scandinavia, setting up kingdoms as they went.
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1
+
2
+
3
+ 3 (three) is a number, numeral, and glyph. It is the natural number following 2 and preceding 4, and is the smallest odd prime number. It has religious or cultural significance in many societies.
4
+
5
+ The use of three lines to denote the number 3 is only natural[clarification needed] and occurred in many writing systems, including some (like Roman and Chinese numerals) that are still in use.
6
+
7
+ In particular, that was also the original representation of 3 in the Brahmin Indians' numerical notation. However, during the Gupta Empire the sign was modified by the addition of a curve on each line. The Nagari rotated the lines clockwise[clarification needed], ended each line with a short downward stroke on the right. In cursive, script the three strokes were eventually connected to form a glyph resembling "3" with an additional stroke at the bottom as "३".
8
+
9
+ The Hindu[clarification needed] numerals spread to the Caliphate in the 9th century. The bottom stroke was dropped around the 10th century in the western parts of the Caliphate, such as the Maghreb and Al-Andalus, when a distinct variant ("Western Arabic") of the digit symbols developed, including modern Western 3. In contrast, the Eastern Arabs retained and enlarged that stroke, rotating the character once more to yield the modern ("Eastern") Arabic digit "٣".[1]
10
+
11
+ In most modern Western typefaces, the "3" glyph, like the other decimal digits, has the height of a capital letter, and sits on the baseline. In typefaces with text figures, on the other hand, the glyph usually has the height of a lowercase letter "x" and a descender: "". In some French text-figure typefaces, though, it has an ascender instead of a descender.[citation needed]
12
+
13
+ A common variant of the digit three has a flat top, similar to the character Ʒ (ezh). This form is sometimes used to obstruct changing a three into an eight with fraudulent intent. It is found on UPC-A barcodes and standard 52-card decks.[citation needed]
14
+
15
+ 3 is:
16
+
17
+ Three is the only prime which is one less than a perfect square. Any other number which is n2 − 1 for some integer n is not prime, since it is (n − 1)(n + 1). This is true for 3 as well (with n = 2), but in this case the smaller factor is 1. If n is greater than 2, both n − 1 and n + 1 are greater than 1 so their product is not prime.
18
+
19
+ A natural number is divisible by three if the sum of its digits in base 10 is divisible by 3. For example, the number 21 is divisible by three (3 times 7) and the sum of its digits is 2 + 1 = 3. Because of this, the reverse of any number that is divisible by three (or indeed, any permutation of its digits) is also divisible by three. For instance, 1368 and its reverse 8631 are both divisible by three (and so are 1386, 3168, 3186, 3618, etc.). See also Divisibility rule. This works in base 10 and in any positional numeral system whose base divided by three leaves a remainder of one (bases 4, 7, 10, etc.).
20
+
21
+ Three of the five Platonic solids have triangular faces – the tetrahedron, the octahedron, and the icosahedron. Also, three of the five Platonic solids have vertices where three faces meet – the tetrahedron, the hexahedron (cube), and the dodecahedron. Furthermore, only three different types of polygons comprise the faces of the five Platonic solids – the triangle, the square, and the pentagon.
22
+
23
+ There are only three distinct 4×4 panmagic squares.
24
+
25
+ According to Pythagoras and the Pythagorean school, the number 3, which they called triad, is the noblest of all digits, as it is the only number to equal the sum of all the terms below it, and the only number whose sum with those below equals the product of them and itself.[2]
26
+
27
+ The trisection of the angle was one of the three famous problems of antiquity.
28
+
29
+ Gauss proved that every integer is the sum of at most 3 triangular numbers.
30
+
31
+ There is some evidence to suggest that early man may have used counting systems which consisted of "One, Two, Three" and thereafter "Many" to describe counting limits. Early peoples had a word to describe the quantities of one, two, and three but any quantity beyond was simply denoted as "Many". This is most likely based on the prevalence of this phenomenon among people in such disparate regions as the deep Amazon and Borneo jungles, where western civilization's explorers have historical records of their first encounters with these indigenous people.[3]
32
+
33
+ Many world religions contain triple deities or concepts of trinity, including:
34
+
35
+ Three is a very significant number in Norse mythology, along with its powers 9 and 27.
36
+
37
+ Three (三, formal writing: 叁, pinyin sān, Cantonese: saam1) is considered a good number in Chinese culture because it sounds like the word "alive" (生 pinyin shēng, Cantonese: saang1), compared to four (四, pinyin: sì, Cantonese: sei1), which sounds like the word "death" (死 pinyin sǐ, Cantonese: sei2).
38
+
39
+ Counting to three is common in situations where a group of people wish to perform an action in synchrony: Now, on the count of three, everybody pull! Assuming the counter is proceeding at a uniform rate, the first two counts are necessary to establish the rate, and the count of "three" is predicted based on the timing of the "one" and "two" before it. Three is likely used instead of some other number because it requires the minimal amount counts while setting a rate.
40
+
41
+ There is another superstition that it is unlucky to take a third light, that is, to be the third person to light a cigarette from the same match or lighter. This superstition is sometimes asserted to have originated among soldiers in the trenches of the First World War when a sniper might see the first light, take aim on the second and fire on the third.
42
+
43
+ The phrase "Third time's the charm" refers to the superstition that after two failures in any endeavor, a third attempt is more likely to succeed. This is also sometimes seen in reverse, as in "third man [to do something, presumably forbidden] gets caught".
44
+
45
+ Luck, especially bad luck, is often said to "come in threes".[19]
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1
+
2
+
3
+ After 1291
4
+
5
+ Northern Crusades (1147–1410)
6
+
7
+ Popular crusades
8
+
9
+ Against Christians
10
+
11
+ Against Ottomans
12
+
13
+ Reconquista (718–1492)
14
+
15
+ The Crusades were a series of religious wars initiated, supported, and sometimes directed by the Latin Church in the medieval period. The term refers especially to the Eastern Mediterranean campaigns in the period between 1096 and 1271 that had the objective of recovering the Holy Land from Islamic rule. The term has also been applied to other church-sanctioned campaigns fought to combat paganism and heresy, to resolve conflict among rival Roman Catholic groups, or to gain political and territorial advantage. The difference between these campaigns and other Christian religious conflicts was that they were considered a penitential exercise that brought forgiveness of sins declared by the church. Historians contest the definition of the term "crusade". Some restrict it to only armed pilgrimages to Jerusalem; others include all Catholic military campaigns with a promise of spiritual benefit; all Catholic holy wars; or those with a characteristic of religious fervour.
16
+
17
+ In 1095, Pope Urban II proclaimed the First Crusade at the Council of Clermont. He encouraged military support for Byzantine Emperor Alexios I against the Seljuk Turks and an armed pilgrimage to Jerusalem. Across all social strata in western Europe there was an enthusiastic popular response. Volunteers took a public vow to join the crusade. Historians now debate the combination of their motivations, which included the prospect of mass ascension into Heaven at Jerusalem, satisfying feudal obligations, opportunities for renown, and economic and political advantage. Initial successes established four Crusader states in the Near East: the County of Edessa; the Principality of Antioch; the Kingdom of Jerusalem; and the County of Tripoli. The crusader presence remained in the region in some form until the city of Acre fell in 1291, leading to the rapid loss of all remaining territory in the Levant. After this, there were no further crusades to recover the Holy Land.
18
+
19
+ Proclaimed a crusade in 1123, the struggle between the Christians and Muslims in the Iberian Peninsula was called the Reconquista by Christians, and only ended in 1492 with the fall of the Muslim Emirate of Granada. From 1147 campaigns in Northern Europe against pagan tribes were considered crusades. In 1199 Pope Innocent III began the practice of proclaiming political crusades against Christian heretics. In the 13th century, crusading was used against the Cathars in Languedoc and against Bosnia; this practice continued against the Waldensians in Savoy and the Hussites in Bohemia in the 15th century and against Protestants in the 16th. From the mid-14th century, crusading rhetoric was used in response to the rise of the Ottoman Empire, only ending in 1699 with the War of the Holy League.
20
+
21
+ In modern historiography, the term "crusade" first referred to military expeditions undertaken by European Christians in the 11th, 12th, and 13th centuries to the Holy Land. The conflicts to which the term is applied has been extended to include other campaigns initiated, supported and sometimes directed by the Roman Catholic Church against pagans, heretics or for alleged religious ends.[1] These differed from other Christian religious wars in that they were considered a penitential exercise, and so earned participants forgiveness for all confessed sins.[2] The term's usage can create a misleading impression of coherence, particularly regarding the early crusades, and the definition is a matter of historiographical debate among contemporary historians.[3][4][5]
22
+
23
+ At the time of the First Crusade, iter, "journey", and peregrinatio, "pilgrimage" were used for the campaign. Crusader terminology remained largely indistinguishable from that of Christian pilgrimage during the 12th century. Only at the end of the century was a specific language of crusading adopted in the form of crucesignatus—"one signed by the cross"—for a crusader. This led to the French croisade—the way of the cross.[3] By the mid 13th century the cross became the major descriptor of the crusades with crux transmarina—"the cross overseas"—used for crusades in the eastern Mediterranean, and crux cismarina—"the cross this side of the sea"—for those in Europe.[6][7] The modern English "crusade" dates to the early 1700s.[8]
24
+
25
+ The Arabic word for struggle or contest, particularly one for the propagation of Islam—jihād—was used for a religious war of Muslims against unbelievers, and it was believed by some Muslims that the Quran and Hadith made this a duty.[9] "Franks" and "Latins" were used by the peoples of the Near East during the crusades for western Europeans, distinguishing them from the Byzantine Christians who were known as "Greeks".[10][11] "Saracen" was used for an Arab Muslim, derived from a Greek and Roman name for the nomadic peoples of the Syro-Arabian desert.[12] Crusader sources used the term "Syrians" to describe Arabic speaking Christians who were members of the Greek Orthodox Church, and "Jacobites" for those who were members of the Syrian Orthodox Church.[13] The Crusader states of Syria and Palestine were known as the "Outremer" from the French outre-mer, or "the land beyond the sea".[14]
26
+
27
+ Christianity was adopted by the Roman Empire in Late Antiquity and Constantinople was founded by the first Christian Roman Emperor, Constantine the Great, in 324. The city developed into the largest in the Christian world, while the Western Roman Empire collapsed at the end of the 5th century. The city and the Eastern Roman Empire are more generally known as Byzantium, the name of the older Greek city it replaced.[15] By the end of the 11th century the period of Islamic Arab territorial expansion had been over for centuries. Its remoteness from focus of Islamic power struggles enabled relative peace and prosperity for the Holy Land in Syria and Palestine. The conflict in the Iberian peninsula was the only location where Muslim-Western European contact was more than minimal.[16]
28
+
29
+ The Byzantine Empire and the Islamic world were long standing centres of wealth, culture and military power. They viewed Western Europe as a backwater that presented little organised threat.[17] The Byzantine Emperor Basil II had extended territorial recovery to its furthest extent in 1025. The Empire's frontiers stretched east to Iran. It controlled Bulgaria, much of southern Italy and suppressed piracy in the Mediterranean Sea. The Empire's relationships with its Islamic neighbours were no more quarrelsome than its relationships with the Slavs or the Western Christians. The Normans in Italy; to the north Pechenegs, Serbs and Cumans; and Seljuk Turks in the east all competed with the Empire and the emperors recruited mercenaries—even on occasions from their enemies—to meet this challenge.[18]
30
+
31
+ After the foundation of the Islamic religion by Muhammad in the 7th century, Muslim Arabs conquered territory from the Indus in the east, and across North Africa and Southern France to the Iberian Peninsula in the West, before political and religious fragmentation halted this expansion. Syria, Egypt, and North Africa were taken from the Byzantine Empire. The emergance of Shia Islam—the belief system that only descendants of Muhammad's cousin and son-in-law, Ali, and daughter, Fatimah, could lawfully be caliph— had led to a split with Sunni Islam on theology, ritual and law. Muslim Iberia was an independent state in modern Spain and Portugal from the 8th century. The Shi'ite Fatimid dynasty ruled North Africa, swathes of Western Asia including Jerusalem, Damascus and parts of the Mediterranean coastline from 969.[19] Total submission to Islam from Jews or Christians was not required. As People of the Book or dhimmi they could continue in their faith on payment of a poll tax. In the Near East a minority Muslim elite ruled over indigenous Christians—Greeks, Armenians, Syrians and Copts.[20]
32
+
33
+ Waves of Turkic migration into the Middle East enjoined Arab and Turkic history from the 9th century. Prisoners from the borderlands of Khurasan and Transoxania were transported to central Islamic lands, converted to Islam and given military training. Known as ghulam or mamluks, it was expected that as slaves they would be more loyal to their masters. In practice it took these Turks only a few decades to progress from being guards, to commanders, governors, dynastic founders and eventually king makers. Examples include the Tulunids in Egypt and Syria (868–905) and the Ikhshidids who followed in Egypt (935–969).[21]
34
+
35
+ The political situation in Western Asia was further changed by later waves of Turkish migration. In particular, the arrival of the Seljuk Turks in the 10th century. Previously a minor ruling clan from Transoxania, they had recently converted to Islam and migrated into Iran to seek their fortune. In the two decades following their arrival they conquered Iran, Iraq and the Near East. The Seljuks and their followers were from the Sunni Islamic tradition which brought them into conflict in Palestine and Syria with the Shi'ite Fatimids.[22] The Seljuks were nomadic, Turkish speaking and occasionally shamanistic, very different to their sedentary, Arabic speaking subjects. This difference and the governance of territory based on political preference, and competition between independent princes rather than geography, weakened power structures.[23] Byzantine Emperor Romanos IV Diogenes attempted confrontation in 1071 to suppress the Seljuks sporadic raiding leading to defeat at the Battle of Manzikert. Historians once considered this a pivotal event but now Manzikert is regarded as only one further step in the expansion of the Great Seljuk Empire.[24]
36
+
37
+ The papacy had declined in power and influence to little more than a localised bishopric by the start of the 11th century. But in the period from the 1050s until the 1080s, under the influence of the Gregorian Reform movement, it became increasingly assertive. Conflict with eastern Christians resulted from the doctrine of papal supremacy. The Eastern church viewed the pope as only one of the five patriarchs of the Church, alongside the Patriarchates of Alexandria, Antioch, Constantinople and Jerusalem. In 1054 differences in custom, creed, and practice spurred Pope Leo IX to send a delegation to the Patriarch of Constantinople, which ended in mutual excommunication and an East–West Schism.[25]
38
+
39
+ The use of violence for communal purposes was not alien to early Christians. The evolution of a Christian theology of war was inevitable when Roman citizenship became linked to Christianity and citizens were required to fight against the Empire's enemies. This was supported by the development of a doctrine of holy war dating from the works of the 4th-century theologian Augustine. Augustine maintained that an aggressive war was sinful, but acknowledged a "just war" could be rationalised if it was proclaimed by a legitimate authority such as a king or bishop, was defensive or for the recovery of lands, and a without an excessive degree of violence.[26][27]
40
+
41
+ Violent acts were commonly used for dispute resolution in Western Europe, and the papacy attempted to mitigate it.[28] Historians, such as Carl Erdmann, thought the Peace and Truce of God movements restricted conflict between Christians from the 10th century; the influence is apparent in Pope Urban II's speeches. But later historians, such as Marcus Bull, assert that the effectiveness was limited and it had died out by the time of the crusades.[29]
42
+
43
+ Pope Alexander II developed a system of recruitment via oaths for military resourcing that Gregory VII extended across Europe. [30] Christian conflict with Muslims on the southern peripheries of Christendom was sponsored by the Church in the 11th century, including the siege of Barbastro and fighting in Sicily[31] In 1074 Gregory VII planned a display of military power to reinforce the principle of papal sovereignty. His vision of a holy war supporting Byzantium against the Seljuks was the first crusade prototype, but lacked support.[32] Theologian Anselm of Lucca took the decisive step towards an authentic crusader ideology, stating that fighting for legitimate purposes could result in the remission of sins.[33]
44
+
45
+ The first crusade was advocated by Urban II at the Council of Clermont in 1095, promising absolution for the participants' sins.[34] An equivalence was created between crusades for the Holy Land and the Reconquista by Calixtus II in 1123. During the period of the Second Crusade Eugenius III was persuaded by the Cistercian abbot, Bernard of Clairvaux, that the German's conquest of the pagan Slavs was also comparable.[35] The 1146 papal bull Divina dispensatione declared pagan conversion was a goal worthy of crusade.[36] Papal protection, penance and salvation for those killed was extended to participants in the suppression of heretical sects in 1179 during the Third Council of the Lateran.[37]
46
+
47
+ Elected pope in 1198, Innocent III reshaped the ideology and practice of crusading. He emphasised crusader oaths and penitence, and clarified that the absolution of sins was a gift from God, rather than a reward for the crusaders' sufferings. Taxation to fund crusading was introduced and donation encouraged.[38][39] In 1199 he was the first pope to deploy the conceptual and legal apparatus developed for crusading to enforce papal rights. With his 1213 bull Quia maior he appealled to all Christians, not just the nobility, offering the possibility of vow redemption without crusading. This set a precedent for trading in spiritual rewards, a practice that scandalised devout Christians and later became one of the causes of the 16th-century Protestant Reformation.[40][41] From the 1220s crusader privileges were regularly granted to those who fought against heretics, schismatics or Christians the papacy considered non-conformist.[42] When Frederick II's army threatened Rome, Gregory IX used crusading terminology. Rome was seen as the Patrimony of Saint Peter, and canon law regarded crusades as defensive wars to protect theoretical Christian territory.[43]
48
+
49
+ Innocent IV rationalised crusading ideology on the basis of the Christians' right to ownership. He acknowledged Muslims' land ownership, but emphasised that this was subject to Christ's authority.[44] In the 16th century the rivalry between Catholic monarchs prevented anti-Protestant crusades but individual military actions were rewarded with crusader privileges, including Irish Catholic rebellions against English Protestant rule and the Spanish Armada's attack on Queen Elizabeth I and England.[45]
50
+
51
+ The First Crusade was an unexpected event for contemporary chroniclers, but historical analysis demonstrates it had its roots in developments earlier in the 11th century. Clerics and laity increasingly recognised Jerusalem as worthy of penitential pilgrimage. In 1071, Jerusalem was captured by the Turkish warlord Atsiz, who seized most of Syria and Palestine as part of the expansion of the Seljuk Turks throughout the Middle East. The Seljuk hold on the city was weak and returning pilgrims reported difficulties and the oppression of Christians. Byzantine desire for military aid converged with increasing willingness of the western nobility to accept papal military direction.[46][47]
52
+
53
+ The desire of Christians for a more effective Church was evident in increased piety. Pilgrimage to the Holy Land expanded after safer routes through Hungary developed from 1000. There was an increasingly articulate piety within the knighthood and the developing devotional and penitential practises of the aristocracy created a fertile ground for crusading appeals.[30] Crusaders' motivations may never be understood. One factor may have been spiritual – a desire for penance through warfare. The historian Georges Duby's explanation was that crusades offered economic advancement and social status for younger, landless sons of the aristocracy. This has been challenged by other academics because it does not account for the wider kinship groups in Germany and Southern France. The anonymous Gesta Francorum talks about the economic attraction of gaining "great booty". This was true to an extent, but the rewards often did not include the seizing of land, as fewer crusaders settled than returned. Another explanation was adventure and an enjoyment of warfare, but the deprivations the crusaders experienced and the costs they incurred weigh against this. One sociological explanation was that crusaders had no choice as they were embedded in extended patronage systems and obliged to follow their feudal lords.[48]
54
+
55
+ From 1092 the status quo in the Middle East disintegrated following the death of the vizier and effective ruler of the Seljuk Empire, Nizam al-Mulk. This was closely followed by the deaths of the Seljuk Sultan Malik-Shah and the Fatimid khalif, Al-Mustansir Billah. The Islamic historian Carole Hillenbrand has described this as analogous to the fall of the Iron Curtain in 1989 with the phrase "familiar political entities gave way to disorientation and disunity".[49] The confusion and division meant the Islamic world disregarded the world beyond; this made it vulnerable to, and surprised by, the First Crusade.[50]
56
+
57
+ In 1095 the Byzantine emperor, Alexios I Komnenos, requested military support from the Council of Piacenza for the fight with the Seljuk Turks. Later that year, at the Council of Clermont, Pope Urban supported this and exhorted war.[51] Thousands of predominantly poor Christians, led by the French priest Peter the Hermit, formed the first response known as the People's Crusade. Passing through Germany they indulged in wide-ranging anti-Jewish activities and massacres. On leaving Byzantine-controlled territory in Anatolia they were annihilated in a Turkish ambush at the Battle of Civetot in October 1096.[52]
58
+
59
+ They were followed by independent military contingents in loose, fluid arrangements based on bonds of lordship, family, ethnicity and language led by members of the high nobility. Foremost were five princes: Count Raymond of Toulouse; two Normans from southern Italy—Bohemond of Taranto and his nephew Tancred—Godfrey of Bouillon; and his brother Baldwin who led a force from Lotharingia, and Germany. They were joined by a northern French army led by Robert Curthose; Count Stephen of Blois; and Count Robert of Flanders. The army may have numbered 100,000 including non-combatant. They travelled east by land and were cautiously welcomed to Byzantium by Alexios late in 1096.[53] He made them promise to return all recovered Byzantine territory and that their first objective should be Nicaea. While the Seljuk Sultan of Rûm, Kilij Arslan, was away resolving a dispute a Frankish siege and Byzantine naval assault captured the city in June 1097. The crusade then embarked on an arduous march across Anatolia, suffering starvation, thirst and disease. The crusaders gained experience in countering the Turkish tactics of employing lightly armoured mounted archers at the Battle of Dorylaeum. They also developed links with local Armenians. Baldwin left with a small force to establish the County of Edessa, the first Crusader state, early in 1098.[54]
60
+
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+ In June 1098 the crusaders gained entry to Antioch after an eight-month siege, massacring most inhabitants, including local Christians. Kerbogha, the Atabeg of Mosul, led a relief force to the city, but Bohemond repulsed him. There was a delay of months while the crusaders decided who would keep the city. This ended on the news that the Fatimid Egyptians had taken Jerusalem from the Seljuks. Despite his promise to Alexios, Bohemond retained Antioch and remained while Raymond led the army along the coast to Jerusalem.[55] Support transported by the Genoese to Jaffa tilted the balance at the siege of Jerusalem, which fell to the Crusaders. They massacred the inhabitants and pillaged the city. Historians believe that contemporary accounts of the numbers killed were exaggerated, but the narrative of massacre reinforced the crusaders' reputation for barbarism.[56] Godfrey secured the Frankish position by defeating an Egyptian force at the Battle of Ascalon.[57]
62
+
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+ Many crusaders now considered their pilgrimage complete and returned to Europe. Only 300 knights and 2,000 infantry remained to defend Palestine. The support of troops from Lorraine enabled Godfrey, over the claims of Raymond, to take the position of Defender of the Holy Sepulchre. A year later the Lorrainers foiled an attempt by Dagobert of Pisa, the papal legate, to make Jerusalem a theocracy on Godfrey's death. Baldwin was chosen as the first Latin king.[58] Bohemond returned to Europe to fight the Byzantines from Italy, but his 1108 expedition ended in failure. Raymond's successors captured the city of Tripoli after his death, with the support of the Genoese. [59] Relations between Edessa and Antioch were variable: they fought together in the crusader defeat at the Battle of Harran; but the Antiocheans claimed suzerainty and attempted to block the return of Count Baldwin—later king of Jerusalem—from his captivity after the battle.[60] The Franks engaged in Near East politics, with Muslims and Christians often fighting on both sides. The expansion of Antioch came to an end in 1119 with a major defeat by the Turks at the Battle of Ager Sanguinis, also known as the Field of Blood.[61]
64
+
65
+ The limited written evidence available from before 1160 indicates the crusade was barely noticed in the Islamic world. This was probably the result of cultural misunderstanding: the Muslims did not recognise the crusaders as religiously motivated warriors intent on conquest and settlement. They assumed this was the latest in a long line of attacks by Byzantine mercenaries. The Islamic world was divided, with rival rulers in Cairo, Damascus, Aleppo and Baghdad. This gave the crusaders an opportunity for consolidation before a pan-Islamic counter-attack.[62]
66
+
67
+ The rise of Imad al-Din Zengi threatened the Franks. He became Atabeg of Mosul in 1127, expanded his control to Aleppo and in 1144 he conquered Edessa. Two years later Pope Eugenius III called for a second crusade. Bernard of Clairvaux spread the message that the loss was the result of sinfulness. Simultaneously, the anti-Semitic preaching of the Cistercian monk, Rudolf, initiated more massacres of Jews in the Rhineland.[63] This was part of a general increase in crusading activity, including in the Iberian peninsula and northern Europe.[64]
68
+
69
+ Zengi was murdered in uncertain circumstances. His elder son Sayf ad-Din succeeded him as atabeg of Mosul while a younger son Nur ad-Din succeeded in Aleppo.[65] Kings Louis VII of France and Conrad III of Germany were the first ruling monarchs to campaign, but the crusade was not a success. Edessa's destruction made its recovery impossible, and the objectives were unclear. The French held the Byzantines responsible for their defeats by the Seljuks in Anatolia, while the Byzantines reiterated claims on any future territorial gains in northern Syria. The crusaders decided to attack Damascus, breaking a long period of cooperation between Jerusalem and the city's Seljuk rulers. Bad luck, poor tactics and a feeble five-day siege of the city led to argument; the barons of Jerusalem withdrew support and the crusaders retreated before Zengi's sons' army. The chronicler William of Tyre related, and modern historians have concurred, that morale fell, hostility to the Byzantines grew and distrust developed between the newly arrived crusaders and those that had made the region their home.[63]
70
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+ Jerusalem demonstrated an increasing interest in expanding into Egyptian territory after the capture of Ascalon in 1153 opened the road south. A year later Nur ad-Din became the first Muslim in the crusading era to unite Aleppo and Damascus.[66] In 1163 King Amalric of Jerusalem initiated a failed invasion of Egypt which prompted Nur ad-Din to move against the Franks and gain a strategic foothold on the Nile. His Kurdish general, Shirkuh, stormed Egypt and only an Egyptian–Jerusalemite alliance forced his return to Syria. Amalric broke the alliance in a series of ferocious attacks and the Egyptians requested military support. Shirkuh was deployed for a second time, accompanied by his nephew, Yusuf ibn Ayyub, who became known by his Arabic honorific Ṣalāḥ ad-Dīn ("the goodness of faith"), which has been westernised as Saladin. Amalric retreated and the Fatimid caliph appointed the Sunni Shirkuh as vizier. Saladin successfully intrigued to become Shirkuh's successor on his death in 1171. Saladin imprisoned the last Fatimids and established a Sunni regime in Egypt.[67]
72
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+ Nur al-Din died in 1174 and Saladin became regent for his 11-year-old son, As-Salih Ismail al-Malik. The prince died seven years later, but Saladin had already seized Damascus and much of Syria from his ward's relatives.[68] Overconfidence led to an initial defeat by the Franks at the Battle of Montgisard, but Saladin established a domain stretching from the Nile to the Euphrates through a decade of politics, coercion and low-level military action.[69] In 1186 a life-threatening illness prompted him to make good on his propaganda as the champion of Islam and intensify the campaign against the Franks.[70] King Guy of Jerusalem responded by raising the largest army that Jerusalem had ever put into the field. This force was lured into inhospitable terrain without water and routed by Saladin's forces at the Battle of Hattin. Numerous Christian nobles were taken prisoner, including Guy. Saladin offered them the option of leaving within 40 days or remaining in peace under Islamic rule. Jerusalem and much of Palestine quickly fell to Saladin.[71]
74
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75
+ Pope Gregory VIII issued a papal bull titled Audita tremendi, that proposed what became known as the Third Crusade. In August 1189, the freed King Guy attempted to recover Acre by surrounding the city and a long stalemate ensued.[72] Travelling overland Holy Roman Emperor Frederick I died crossing the Saleph River in Cilicia and only a few of his men reached their destination. King Richard I of England travelled by sea. Philip II of France was the first king to arrive at the siege.[73] Richard I conquered Cyprus in transit in response to his sister and his fiancée being take prisoner by the Cypriot ruler, Isaac Komnenos.[74] A year later Richard would facilitate the sale of the island to King Guy for 40,000 bezants as part of the settlement replacing Guy as king of Jerusalem with Conrad of Montferrat.[75]
76
+
77
+ The arrival of the French and English turned the tide in the conflict, and the Muslim garrison of Acre surrendered. Philip considered his vow fulfilled and returned to France, leaving most of his forces behind. Richard travelled south along the Mediterranean coast and recaptured Jaffa. Twice he advanced to within a day's march of Jerusalem, but lacked the resources to capture and defend the city. A negotiated three-year truce allowed Frankish access to Jerusalem. This was the end of Richard's crusading career and damaged Frankish morale.[76] The Crusader states survived, confined to a narrow coastal strip.[77] Emperor Frederick I's successor, Henry VI, announced a new crusade without papal encouragement in 1195. Henry died before departing on the crusade, but the arrival of the German crusaders prompted Saladin's brother, Al-Adil I to sign a five-year truce in 1198.[78]
78
+
79
+ In 1198 the recently elected Pope Innocent III announced a new crusade, organised by three Frenchmen: Theobald of Champagne; Louis of Blois; and Baldwin of Flanders. The Italian Boniface of Montferrat replaced Theobald on the latter's premature death, as the new commander of the campaign. They contracted with the Republic of Venice for the transportation of 30,000 crusaders at a cost of 85,000 marks. However, many choose other embarkation ports and only around 15,000 arrived at Venice. Unable to fully pay the Venetians they accepted two offers. The Doge of Venice Enrico Dandolo proposed that Venice would be repaid with the profits of future conquests beginning with the seizure of the Christian city of Zara. Secondly, the exiled Byzantine prince Alexios Angelos offered 10,000 troops, 200,000 marks and the reunion of the Greek Church with Rome if they toppled his uncle Emperor Alexios III.[79]
80
+
81
+ Innocent III excommunicated the crusaders for their capture of Zara, but quickly absolved the French. The crusade entered Constantinople, Alexios III fled and was replaced by his nephew. The Greeks resisted the imposition of Alexios IV and harried the crusaders, so he encouraged the crusade to support him until he could fulfil his commitments. This situation ended in a violent anti-Latin revolt and the assassination of Alexios IV. Without ships, supplies or food the crusaders had little option than to take by force what Alexios had promised. The Sack of Constantinople involved three days pillaging churches and killing much of the Greek Orthodox Christian populance.[80] While not unusual behaviour for the time, contemporaries such as Innocent III and Ali ibn al-Athir saw it as an atrocity against centuries of classical and Christian civilisation.[81]
82
+
83
+ A council of six Venetians and six Franks partitioned the territorial gains, establishing a Latin Empire. Baldwin became Emperor of seven-eights of Constantinople, Thrace, northwest Anatolia and the Aegean Islands. Venice gained a maritime domain including the remaining portion of the city. Boniface received Thessalonika, and his conquest of Attica and Boeotia formed the Duchy of Athens. His vassals, William of Champlitte and Geoffrey of Villehardouin, conquered Morea, establishing the Principality of Achaea. Both Baldwin and Boniface died fighting the Bulgarians, leading the papal legate to release the crusaders from their obligations.[82][83] As many as a fifth of the crusaders continued to Palestine via other routes, including a large Flemish fleet. Joining King Aimery on campaign they forced al-Adil into a six-year truce.[84]
84
+
85
+ The Latin states established were a fragile patchwork of petty realms threatened by Byzantine successor states—the Despotate of Epirus, the Empire of Nicaea and the Empire of Trebizond. Thessaloniki fell to Epirus in 1224, and Constantinople to Nicaea in 1261. Achaea and Athens survived under the French after the Treaty of Viterbo.[85][86] The Venetians endured a long-standing conflict with the Ottoman Empire until the final possessions were lost in the Seventh Ottoman–Venetian War in the 18th century. This period of Greek history is known as the Frankokratia or Latinokratia ("Frankish or Latin rule") and designates a period when western European Catholics ruled Orthodox Byzantine Greeks.[87]
86
+
87
+ There were repeated popular outbursts of ecstatic piety in 13th-century Western Europe such as the Children's Crusade of 1212, when large groups of young adults and children gathered spontaneously in the belief that their innocence would lead to success where others had failed. Few, if any, journeyed to the eastern Mediterranean.[40] Crusading did not resume until 1217. There was no immediate threat and a number of treaties had to expire first. Little was achieved by a Fifth Crusade, primarily raised from Hungary, Germany, Flanders and led by King Andrew II of Hungary and Leopold VI, Duke of Austria. The crusaders attacked Egypt to break the Muslim hold of Jerusalem. Egypt was isolated from the other Islamic power centres, it would be easier to defend and was self-sufficient in food. Damietta was captured but then returned and an eight-year truce agreed after the Franks advancing into Egypt surrendered.[88]
88
+
89
+ Holy Roman Emperor Frederick II had frequently postponed fulfilling his crusading commitments before he acquired the Kingdom of Jerusalem through marriage in 1225. In 1227 he embarked on crusade, but was forced to abandon it due to illness. This prompted his excommunication by Pope Gregory IX. Despite this Frederick launched a campaign of forceful negotiation that won the Franks most of Jerusalem, a strip of territory linking the city to Acre and an alliance with Al-Kamil, Sultan of Egypt. When the Pope attacked Frederick's Italian possessions he returned to defend them.[89] The kingdom could no longer rely on Frederick's resources and was left dependent on Ayyubid division, the military orders and western aid for survival.[90] The popes' conflict with Frederick left the responsibility for crusading to secular, rather than papal, leadership. The Barons' Crusade was led by King Theobald I of Navarre and when he returned home, by the king of England's brother, Richard of Cornwall. The Franks followed Frederick's tactics of forceful diplomacy and playing rival factions off against each other when Sultan Al-Kamil died and his family fell into disputes over the succession in Egypt and Syria.[91]
90
+
91
+ The Mongols provided a new military threat to the Christian and Islamic worlds, sweeping west through southern Russia, Poland and Hungary; defeating the Seljuks and threatening the Crusader states. Although predominantly pagan, some Mongols were Nestorian Christians. This gave the papacy hope they might become allies. But when Pope Innocent IV wrote to the Mongols to question their attacks on Christians they replied demanding his total submission.[92] The Mongols displaced a central Asian Turkish people, the Khwarazmian, providing Al-Kamil's son As-Salah with useful allies. The Khwarazmians captured Jerusalem and savagely sacked it. An Egyptian–Khwarazmian army then annihilated a Frankish–Damascene army at the Battle of La Forbie. This was the last time the Franks had the resources to raise a field army in Palestine. As-Salah conquered almost all of the crusaders' mainland territories, confining them to a few coastal towns.[93][94]
92
+
93
+ The devout French king, Louis IX, and his brother, Charles I of Anjou, dominated 13th-century politics in the eastern Mediterranean. In 1249 Louis led a crusade attacking Egypt and was defeated at the Battle of Al Mansurah and the crusaders were captured as they retreated. Louis and his nobles were ransomed, other prisoners were given a choice of conversion to Islam or beheading. A ten-year truce was established and Louis remained in Syria until 1254 consolidating the Frankish position. In Egypt a power struggle developed between the Mamluks and the Ayyubid rulers. This led to one of the Mamluk leaders, Qutuz, seizing the sultanate in 1259 and uniting with another Mamluk faction led by Baibars. The Mamluks defeated the Mongols at Ain Jalut before gaining control of Damascus and Aleppo. Qutuz was assassinated and Baibars assumed control.[95][96]
94
+
95
+ Division in the crusader states led to conflicts such as the War of Saint Sabas. Venice drove the Genoese from Acre to Tyre where they continued trading with the Egyptians.[97] In 1270 Charles turned Louis's new crusade to his advantage by persuading him to attack Tunis. Their army was devastated by disease, and Louis died at Tunis. Prince Edward, the future king of England, and a small retinue arrived too late for the conflict but continued to the Holy Land. Edward survived an assassination attempt, negotiated a ten-year truce, and then returned to manage his affairs in England. This ended the last significant crusading effort in the eastern Mediterranean.[98] The mainland crusader states were finally extinguished with the fall of Tripoli in 1289 and Acre in 1291.[99]
96
+
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+ The causes of the decline in crusading and the failure of the crusader states are multi-faceted. The nature of crusades was unsuited to the defence of the Holy Land. Crusaders were on a personal pilgrimage and usually returned when it was completed. Although the ideology of crusading changed over time, crusades continued to be conducted without centralised leadership by short-lived armies led by independently minded potentates, but the crusader states needed large standing armies. Religious fervour was difficult to direct and control even though it enabled significant feats of military endeavour. Political and religious conflict in Europe combined with failed harvests reduced Europe's interest in Jerusalem. The distances involved made the mounting of crusades and the maintenance of communications difficult. It enabled the Islamic world, under the charismatic leadership of Zengi, Nur al-Din, Saladin, the ruthless Baibars and others, to use the logistical advantages of proximity.[100]
98
+
99
+ After the First Crusade most of the crusaders considered their personal pilgrimage complete and returned to Europe.[57] Modern research indicates that Muslim and indigenous Christian populations were less integrated than previously thought. Palestinian Christians lived around Jerusalem and in an arc stretching from Jericho and the Jordan to Hebron in the south.[101] Archaeological research on Byzantine churches and Ottoman census records from the 16th century demonstrate that Greek Orthodox communities survived centuries after the fall of the Crusader states. Maronites were concentrated in Tripoli, the Jacobites in Antioch and Edessa. Armenians also lived in the north but communities existed in all major towns. Central areas had a Muslim majority population, predominantly Sunni but with Shi'ite communities in Galilee. Druze Muslims lived in the mountains of Tripoli. The Jewish population resided in coastal towns and some Galilean villages.[102][103] The Frankish population of the Kingdom of Jerusalem was concentrated in three major cities. By the 13th century the population of Acre probably exceeded 60,000, then came Tyre and the capital itself was the smallest of the three with a population somewhere between 20,000 and 30,000.[104] The Latin population of the region peaked at c250,000 with Jerusalem's population numbering c120,000 and the combined total in Tripoli, Antioch and Edessa being broadly comparable.[105] In context, Josiah Russell roughly estimates the population of what he calls "Islamic territory" as 12.5 million in 1000 with the European areas that provided crusaders having a population of 23.7 million. He estimates that by 1200 that these figures had risen to 13.7 million in Islamic territory while the Crusaders' home countries population was 35.6 million. Russell acknowledges that much of Anatolia was Christian or under the Byzantines and "Islamic" areas such as Mosul and Baghdad had significant Christian populations.[106]
100
+
101
+ The Outremer was a frontier society in which a Frankish elite ruled over of a native population related to the neighbouring communities, many of whom were hostile to the Franks.[107] It was politically and legally stratified with self-governing ethnic communities. Relations between communities were controlled by the Franks.[108] The basic division in society was between Frank and non-Frank, and not between Christian and Muslim. All Franks were considered free men while the native peoples lived like western serfs. The Franks imposed officials in the military, legal and administrative systems using the law and lordships to control the natives. Few Franks could speak more than basic Arabic. Dragomans—interpreters—and ruʾasāʾ—village headmen—were used as mediators. Civil disputes and minor criminality were administered by the native communities, but major offences and those involving Franks were dealt by the Frankish cour des bourgeois. The key differentiator in status and economic position was between urban and rural dwellers. Indigenous Christians could gain higher status and acquire wealth through commerce and industry in towns but few Muslims lived in urban areas except servants.[109]
102
+
103
+ The Crusader States presented an obstacle to Muslim trade with the west by sea and the land routes from Mesopotamia and Syria to the urban economies of the Nile. However, despite this commerce continued, coastal cities remained maritime outlets for the Islamic hinterland, Eastern wares were exported to Europe in unprecedented volumes. Byzantine-Muslim mercantile growth in the 12th and 13th  centuries may have occurred anyway. Western Europe’s population, wealth and the demand for sophisticated Eastern products was booming but it is likely that the Crusades hastened the developments. European fleets expanded, better ships were built, navigation improved and fare paying pilgrims subsidised many voyages. The mainly native agricultural production flourished before the fall of the First Kingdom in 1187, but was negligible afterwards. Italian, Provençal and Catalan merchants monopolised shipping, imports, exports, transportation and banking while the income of the Franks was based on income from estates, market tolls and taxation.[110] Production centred in Antioch, Tripoli, Tyre and Beirut. The Franks exported textiles, glass dyestuffs, olives, wine, sesame oil, sugar and prized Silk and imported clothing and finished goods.[111] The indigenous monetised economic system was adopted with northern Italian and southern French silver European coins, Frankish copper coins minted in Arabic and Byzantine styles, local silver and gold dirhams and dinars. After 1124, Egyptian dinars were copied creating Jerusalem's gold bezant. Following the collapse of the First Kingdom in 1187, trade rather than agriculture increasingly dominated the economy and western coins dominated the coinage and despite some local minting of silver pennies and coppers there is little evidence of systematic attempts to create a unified local currency.[112]
104
+
105
+ During the near constant warfare in the early decades of the 12th century, the king of Jerusalem's foremost role was leader of the feudal host. They rewarded their followers' loyalty with city incomes rarely granting land and when holdings became vacant, due to the conflict’s high mortality rate this reverted to the crown. The result was that the royal domain of the first five rulers was greater than the combined holdings of the nobility. This gave the rulers of Jerusalem greater internal power than comparative western monarchs but without the necessary administrative machinery to govern a large realm.[113] Baronial dynasties evolved in the second quarter of the century often acting as autonomous rulers. Royal powers were abrogated and effectively governance undertaken locally within the feudatories. Central control that remained was exercised through the Haute Cour or High Court. This was meetings between the king and his tenants in chief. The duty of the vassal to give counsel became a privilege until the legitimacy of the monarch depended on the agreement of the court.[114] The barons have been poorly regarded by both contemporary and modern commentators who note their superficial rhetoric, pedantry and spurious legal justification for political action.[115]
106
+
107
+ The High Court consisted of the great barons and the king's direct vassals with a quorum of the king and three tenants in chief. The 1162 Assise sur la ligece expanded membership to all the 600+ Franks who paid homage directly to the king. They were joined by the heads of the military orders before the end of the 12th century and the Italian communes in the 13th century.[116] Before the defeat at Hattin in 1187 the laws developed were documented as Assises in Letters of the Holy Sepulchre.[117] The entire body of written law was lost in the fall of Jerusalem leaving a legal system largely based on the custom and memory of the lost legislation. Philip of Novara wrote We know [the laws] rather poorly, for they are known by hearsay and usage...and we think an assize is something we have seen as an assize...in the kingdom of Jerusalem [the barons] made much better use of the laws and acted on them more surely before the land was lost. A myth was created of an idyllic early 12th century legal system that the barons used to reinterpret the Assise sur la ligece that Almalric I intended to strengthen the crown to rather than constrain the monarch’s ability to confiscate feudal fiefs without trial. When the rural fiefs were lost the barons became an urban mercantile class whose knowledge of the law was a valuable skill and career path to higher status.[118] The leaders of the Third Crusade considered the monarchy of Jerusalem of secondary importance. They decided on the grants of land and even granted the throne itself in 1190 and 1192, to Conrad of Montferrat and Henry II, Count of Champagne.[119] Emperor Frederick II married Queen Isabella in 1225 and claimed the throne from her father, the King Regent—John of Brienne. In 1228 Isabella II died after giving birth to a son, Conrad, who through his mother was now legally king of Jerusalem and Frederick's heir.[89] Frederick II left the Holy Land to defend his Italian and German lands beginning a period of absent monarchs from 1225 until 1254. In contrast to Western monarchies with powerful, with centralised bureaucracies government in Jerusalem developed in the opposite direction. Jerusalem’s royalty had title but little power.[120] Magnates fought for regency control with an Italian army led by Frederick's viceroy Richard Filangieri in the War of the Lombards. Tyre, the Hospitallers, the Teutonic Knights and Pisa supported Filangieri. In opposition were the Ibelins, Acre, the Templars and Genoa. For twelve years the rebels held a surrogate parliament in Acre before prevailing in 1242, leading toy a succession of Ibelin and Cypriot regents .[121][122] Centralised government collapsed and the nobility, military orders and Italian communes took the lead. Three Cypriot Lusignan kings succeeded without the resources to recover the lost territory. The title of king was sold to Charles of Anjou who gained power for a short while but never visited the kingdom. [123]
108
+
109
+ The early crusaders filled ecclesiastical positions left vacant by the Orthodox church and replaced Orthodox bishops with Latin clerics. The Greek Orthodox monks of the Holy Sepulchre were expelled but recalled when the miracle of Easter Fire failed in their absence. Armenians, Copts, Jacobites, Nestorians and Maronites were considered autonomous, retaining their own bishops.[124] Assimilation was prevented by discriminatory laws for Jews and Muslims and an absence of effort by the Franks. Muslims were banned from living in Jerusalem and sexual relationships between Muslims and Christians was punished by mutilation. [125]
110
+
111
+ Largely based in the ports of Acre, Tyre, Tripoli and Sidon, Italian, Provençal and Catalan communes had distinct cultural characteristics and exerted significant political power. Separate from the Frankish nobles or burgesses, the communes were autonomous political entities closely linked to their hometowns. They monopolised foreign trade and almost all banking and shipping and aggressively extended trade privileges. Despite all efforts, the ports were unable supersede Alexandria and Constantinople as the primary regional commercia centres but the communes did compete with the monarchs and each other for economic advantage. Power derived from the support of the communards' native cities rather than their number, which never reached more than hundreds. By the middle of the 13th century, the rulers of the communes were barely recognised crusader authority and divided Acre into several fortified miniature republics.[126][127]
112
+
113
+ There were few cultural innovations in the Outremer beyond the establishment of the military orders and the development of tactics and military architecture.[128] John of Ibelin records in around 1170 that the military force of the kingdom of Jerusalem was based on a feudal host of about 647 to 675 heavily armoured knights. Each knight would also provide his own armed retainers. Non-noble light cavalry and infantry were known as serjants and these numbered around 5,025. These numbers were augmented by mercenaries such as the Turcopoles recruited from among the natives. [129] Joshua Prawer estimated that the military orders matched this force in number giving an estimated military strength of 1,200 knights and 10,000 serjants. This was sufficient for territorial gains, but fewer than the required to maintain military domination. This defensive problem was that putting an army into the field required draining castles and cities of every able-bodied fighting man. In the case of a defeat such as at Hattin, no one remained to resist the invaders. The Franks adopted delaying tactics when faced with a superior invading Muslim force, avoiding direct confrontation, retreating to strongholds and waiting for the Muslim army to disperse. Muslim armies were incohesive and seldom campaigned beyond a period between sowing and harvest. It took generations before the Muslims identified that in order to conquer the Crusader states they needed to destroy the Frankish fortresses. This strategic change forced the crusaders away from focussing on the gaining and holding territory but rather on attacking and destroying Egypt, neutralising this regional challenge and gaining the time to improve the kingdom's demographic weaknesses.[130]
114
+
115
+ The disintegration of the Caliphate of Córdoba in southern Spain created the opportunity for the Reconquista, beginning in 1031. The Christian realms had no common identity or shared history based on tribe or ethnicity. As a result, León, Navarre and Catalonia united and divided several times during the 11th and 12th centuries. Although small, all developed an aristocratic military technique.[131] By the time of the Second Crusade the three kingdoms were powerful enough to conquer Islamic territory—Castile, Aragon and Portugal.[132] In 1212 the Spanish were victorious at the Battle of Las Navas de Tolosa with the support of 70,000 foreign combatants who responded to the preaching of Innocent III. Many foreigners deserted because of the tolerance the Spanish demonstrated for the defeated Muslims. For the Spanish, the Reconquista was a war of domination rather than a war of extermination.[133] This contrasted with the treatment of the Christians formerly living under Muslim rule, the Mozarabs. The Roman Rite was relentlessly imposed on them, and the native Christians were absorbed into mainstream Catholicism.[101] Al-Andalus, Islamic Spain, was completely suppressed in 1492 when the Emirate of Granada surrendered. At this point the remaining Muslim and Jewish inhabitants were expelled from the peninsula.[134]
116
+
117
+ There were modest efforts to suppress a dualistic Christian sect called the Cathars in southern France around 1180.[37] After a thirty-year delay Innocent III proclaimed the Albigensian Crusade, named after the city of Albi, one of the centres of Catharism.[135] This proved that it was more effective waging a war against the heretics' supporters than the heretics themselves. Tolerant feudal lords had their lands confiscated and titles forfeited. In 1212 pressure was exerted on the city of Milan for tolerating Catharism.[136] Two Hungarian invasions of Bosnia, the home of a legendary Cathar anti-pope, were proclaimed crusades in 1234 and 1241. A crusade forced the Stedinger peasants of north-western Germany to pay tithes in 1234.[137] The historian Norman Housley notes the connection between heterodoxy and anti-papalism in Italy. Indulgences were offered to anti-heretical groups such as the Militia of Jesus Christ and the Order of the Blessed Virgin Mary.[138] Anti-Christian crusading declined in the 15th century, the exceptions were the six failed crusades against the religiously radical Hussites in Bohemia and attacks on the Waldensians in Savoy.[45]
118
+
119
+ The Albigensian Crusades established a precedent for popes and the Inquisition to claim their Christian opponents were heretics.[139][140] When Frederick threatened to take Rome in 1240, Gregory IX used crusading terminology to raise support. On Frederick II's death the focus moved to Sicily. Until his death the regent, Markward von Annweiler, faced a crusade by Innocent III. In 1263, Pope Urban IV offered crusading indulgences to Charles of Anjou in return for Sicily's conquest. But, these wars had no clear objectives or limitations making them unsuitable for crusading.[43] The 1281 election of a French pope, Martin IV, brought the power of the papacy behind Charles. Charles's preparations for a crusade against Constantinople were foiled by the Byzantine Emperor Michael VIII Palaiologos, who instigated an uprising called the Sicilian Vespers. Instead, Peter III of Aragon was proclaimed king of Sicily, despite his excommunication and an unsuccessful Aragonese Crusade.[141] Political crusading continued against Venice over Ferrara; Louis IV, King of Germany when he marched to Rome for his imperial coronation; and the free companies of mercenaries.[142]
120
+
121
+ In 1147 Bernard of Clairvaux persuaded Pope Eugenius III that the Germans' and Danes' conflict with the pagan Wends was a holy war analogous to the Reconquista; he urged a crusade until all heathens were baptised or killed. The new crusaders' motivation was primarily economic: the acquisition of new arable lands and serfs; the control of Baltic trade routes; and the abolishment of the Novgorodian merchants' monopoly of the fur trade.[143] From the early 13th century the military orders provided garrisons in the Baltic and defended the German commercial centre, Riga. The Livonian Brothers of the Sword and the Order of Dobrzyń were established by local bishops. The Sword Brothers were notorious for cruelty to pagans and converts alike. The Teutonic Knights were founded during the 1190s in Palestine, but their strong links to Germany diverted efforts from the Holy Land to the Baltic. Between 1229 and 1290, the Teutonic Knights absorbed both the Brothers of the Sword and the Order of Dobrzyń, subjugated most of the Baltic tribes and established a ruthless and exploitative monastic state.[144][145] The Knights invited foreign nobility to join their regular Reisen, or raids, against the last unconquered Baltic people, the Lithuanians. These were fashionable events of chivalric entertainment among young aristocrats. Jogaila, Grand Prince of Lithuania, converted to Catholicism and married Queen Jadwiga of Poland resulting in a united Polish–Lithuanian army routing the Knights at Tannenberg in 1410. The Knights' state survived, from 1466 under Polish suzerainty. Prussia was transformed into a secular duchy in 1525, and Livonia in 1562.[146]
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+ The Seljuk Sultanate of Rum fragmented in the late 13th century. The Ottoman Turks, located in north-eastern Anatolia, took advantage of a Byzantine civil war of 1341–1347 and established a strong presence in Europe. They captured the Byzantine fortress at Gallipoli in 1354 and defeated the Serbians at the Battle of Kosovo in 1389, winning control of the Balkans from the Danube to the Gulf of Corinth. This was further confirmed by victory over French crusaders and King Sigismund of Hungary at the Battle of Nicopolis in 1396. Sultan Murad II destroyed a large crusading Serbian and Hungarian force at Varna in 1444 and four years later defeated the Hungarians at Kosovo again.[147][148]
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+ After the fall of Constantinople in 1453 the crusading response was largely symbolic. One example was Duke Phillip of Burgundy's 1454 promotion of a crusade, that never materialised, at the Feast of the Pheasant.[149] The 16th century saw growing rapprochement. The Habsburgs, French, Spanish and Venetians all signed treaties with the Ottomans. King Francis I of France sought allies from all quarters, including from German Protestant princes and Sultan Suleiman the Magnificent.[150] Crusading became chiefly a financial exercise with precedence given to the commercial and political aspects. As the military threat presented by the Turks diminished, anti-Ottoman crusading became obsolete with the Holy League in 1699.[151]
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+ The crusaders' propensity to follow the customs of their Western European homelands meant that there were few innovations developed in the crusader states. Three notable exceptions to this were the military orders, warfare and fortifications.[152] The Knights Hospitaller, formally the Order of Knights of the Hospital of Saint John of Jerusalem, had a medical function in Jerusalem before the First Crusade. The order later adding a martial element and became a much larger military order.[153] In this way knighthood entered the previously monastic and ecclesiastical sphere.[154] The Templars, formally the Poor Fellow-Soldiers of Christ and the Temple of Solomon were founded around 1119 by a small band of knights who dedicated themselves to protecting pilgrims en route to Jerusalem.[155] King Baldwin II granted the order the Al-Aqsa Mosque in 1129 they were formally recognised by the papacy at the 1129 Council of Troyes. Military orders like the Knights Hospitaller and Knights Templar provided Latin Christendom's first professional armies in support of the Kingdom of Jerusalem and the other crusader states.[156]
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+ The Hospitallers and the Templars became supranational organisations as papal support led to rich donations of land and revenue across Europe. This, in turn, led to a steady flow of new recruits and the wealth to maintain multiple fortifications in the crusader states. In time, they developed into autonomous powers in the region.[157] After the fall of Acre the Hospitallers relocated to Cyprus, then ruled Rhodes until the island was taken by the Ottomans in 1522, and Malta until Napoleon captured the island in 1798. The Sovereign Military Order of Malta continues in existence to the present-day.[158] King Philip IV of France probably had financial and political reasons to oppose the Knights Templar, which led to him exerting pressure on Pope Clement V. The Pope responded in 1312 with a series of papal bulls including Vox in excelso and Ad providam that dissolved the order on the alleged and probably false grounds of sodomy, magic and heresy.[159]
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+ According to the historian Joshua Prawer no major European poet, theologian, scholar or historian settled in the crusader states. Some went on pilgrimage, and this is seen in new imagery and ideas in western poetry. Although they did not migrate east themselves, their output often encouraged others to journey there on pilgrimage.[160]
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+ Historians consider the crusader military architecture of the Middle East to demonstrate a synthesis of the European, Byzantine and Muslim traditions and to be the most original and impressive artistic achievement of the crusades. Castles were a tangible symbol of the dominance of a Latin Christian minority over a largely hostile majority population. They also acted as centres of administration.[161] Modern historiography rejects the 19th-century consensus that Westerners learnt the basis of military architecture from the Near East, as Europe had already experienced rapid development in defensive technology before the First Crusade. Direct contact with Arab fortifications originally constructed by the Byzantines did influence developments in the east, but the lack of documentary evidence means that it remains difficult to differentiate between the importance of this design culture and the constraints of situation. The latter led to the inclusion of oriental design features such as large water reservoirs and the exclusion of occidental features such as moats.[162]
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+ Typically, crusader church design was in the French Romanesque style. This can be seen in the 12th-century rebuilding of the Holy Sepulchre. It retained some of the Byzantine details, but new arches and chapels were built to northern French, Aquitanian and Provençal patterns. There is little trace of any surviving indigenous influence in sculpture, although in the Holy Sepulchre the column capitals of the south facade follow classical Syrian patterns.[163]
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+ In contrast to architecture and sculpture, it is in the area of visual culture that the assimilated nature of the society was demonstrated. Throughout the 12th and 13th centuries the influence of indigenous artists was demonstrated in the decoration of shrines, paintings and the production of illuminated manuscripts. Frankish practitioners borrowed methods from the Byzantines and indigenous artists and iconographical practice leading to a cultural synthesis, illustrated by the Church of the Nativity. Wall mosaics were unknown in the west but in widespread use in the crusader states. Whether this was by indigenous craftsmen or learnt by Frankish ones is unknown, but a distinctive original artistic style evolved.[164]
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+ Manuscripts were produced and illustrated in workshops housing Italian, French, English and local craftsmen leading to a cross-fertilisation of ideas and techniques. An example of this is the Melisende Psalter, created by several hands in a workshop attached to the Holy Sepulchre. This style could have both reflected and influenced the taste of patrons of the arts. But what is seen is an increase in stylised, Byzantine-influenced content. This extended to the production of icons, unknown at the time to the Franks, sometimes in a Frankish style and even of western saints. This is seen as the origin of Italian panel painting.[165] While it is difficult to track illumination of manuscripts and castle design back to their origins, textual sources are simpler. The translations made in Antioch are notable, but they are considered of secondary importance to the works emanating from Muslim Spain and from the hybrid culture of Sicily.[166]
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+ Until the requirement was abolished by Innocent III married men needed to obtain their wives' consent before taking the cross, which was not always readily forthcoming. Muslim and Byzantine observers viewed with disdain the many women who joined the armed pilgrimages, including female fighters. Western chroniclers indicated that female crusaders were wives, merchants, servants and sex workers. Attempts were made to control the women's behaviour in ordinances of 1147 and 1190. Aristocratic women had a significant impact: Ida of Formbach-Ratelnberg led her own force in 1101; Eleanor of Aquitaine conducted her own political strategy; and Margaret of Provence negotiated her husband Louis IX's ransom with an opposing woman—the Egyptian sultana Shajar al-Durr. Misogyny meant that there was male disapproval; chroniclers tell of immorality and Jerome of Prague blamed the failure of the Second Crusade on the presence of women. Even though they often promoted crusading, preachers would typecast them as obstructing recruitment, despite their donations, legacies and vow redemptions. The wives of crusaders shared their plenary indulgences.[167]
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+ Despite the common misconception concerning deliberate, widespread pillaging the First Crusade, from its onset, was never intended as a greed-ridden escapade; in fact, Pope Urban II, in his speech at the Council of Clermont, urged his acolytes to "rent their lands and collect money for their expenses" rather than relying on the acquired treasures.[169] The Pope, self-depicted as such in his speech, was not driven by avarice and did not condone the pillaging of Christian and Muslim provinces that were to be encountered throughout the journey. That said, plundering did occur.[170] The mob led by Peter the Hermit ransacked Christian villages in Hungary and Greece, and robbed its inhabitants.[171] Numerous Jewish populaces were robbed and, for some, murdered. But the spoils belonging to the Muslims, significantly greater than those of the two aforementioned religious groups, supplied the crusading forces with an abundance of precious goods, and deposits of gold and silver.[170]
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+ For most participants, however, crusading was ruinously expensive and required extensive financial preparation. The prudent crusader, cognizant of both his militant and at-home obligations, endowed his family, if possible, with the requisite coinage before setting out to crusade, taking with him "great bags and chests of money."[172] That said, Fred Cazel noted: "Few crusaders had sufficient cash both to pay their obligations at home and to support themselves decently on a crusade."[173] Those adamant crusaders, unwilling to forego battle, sold their estates, stocks, and valuables to subsidize the journey ahead; and some even borrowed significant funds from kings and princes, from bishops and monasteries, from merchants and craftsmen, from whoever could sustain lending the ample amount.[170]
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+ The Crusades created national mythologies, tales of heroism, a few place names, and developed Europe's political topology.[174] Historical parallelism and the tradition of drawing inspiration from the Middle Ages have become keystones of political Islam encouraging ideas of a modern jihad and a centuries-long struggle against Christian states, while secular Arab nationalism highlights the role of western imperialism.[175] Modern Muslim thinkers, politicians and historians have drawn parallels between the crusades and political developments such as the establishment of Israel in 1948.[176] Right-wing circles in the western world have drawn opposing parallels, considering Christianity to be under an Islamic religious and demographic threat that is analogous to the situation at the time of the crusades. Crusader symbols and anti-Islamic rhetoric are presented as an appropriate response, even if only for propaganda purposes. These symbols and rhetoric are used to provide a religious justification and inspiration for a struggle against a religious enemy.[177]
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+ Originally, medieval understanding of the crusades was narrowly focussed on a limited set of interrelated texts, most notably Gesta Francorum which possibly dates from as early as 1099. The Gesta was reworked by Robert of Rheims who created a papalist, northern French template for later works. These all demonstrated a degree of martial advocacy that attributed both success and failure to God's will.[178] This clerical view was soon challenged by vernacular adventure stories based on the work of Albert of Aachen. William of Tyre expanded on Albert's writing in his Historia. Completed by 1184, William's work describes the warrior state that Outremer had become through the tensions between divine providence and humankind.[179] Medieval crusade historiography remained more interested in presenting moralistic lessons than information, extolling the crusades as a moral exemplar and a cultural norm.[180]
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+ Attitudes toward the crusades during the Reformation were shaped by confessional debates and the Ottoman expansion. The Protestant martyrologist John Foxe in his History of the Turks (1566) blamed the sins of the Catholic Church for the failure of the crusades. He also condemned the use of crusades against those he considered had maintained the faith, such as the Albigensians and Waldensians. The Lutheran scholar Matthew Dresser (1536–1607) extended this view; the crusaders were lauded for their faith but Urban II's motivation was seen as part of his conflict with Emperor Henry IV. On this view, the crusade was flawed, and the idea of restoring the physical holy places was "detestable superstition".[181] The French Catholic lawyer Étienne Pasquier (1529–1615) was one of the first to number the crusades; he suggested there were six. His work highlights the failures of the crusades and the damage that religious conflict had inflicted on France and the church; it lists victims of papal aggression, sale of indulgences, church abuses, corruption, and conflicts at home.[182]
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+ Age of Enlightenment philosopher-historians such as David Hume, Voltaire and Edward Gibbon used crusading as a conceptual tool to critique religion, civilisation and cultural mores. For them the positives effects of crusading, such as the increasing liberty that municipalities were able to purchase from feudal lords, were only by-products. This view was then criticised in the 19th century by crusade enthusiasts as being unnecessarily hostile to, and ignorant of, the crusades.[183] Alternatively, Claude Fleury and Gottfried Wilhelm Leibniz proposed that the crusades were one stage in the improvement of European civilisation; that paradigm was further developed by the Rationalists.[184]
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+ The idea that the crusades were an important part of national history and identity continued to evolve. In scholarly literature, the term "holy war" was replaced by the neutral German kreuzzug and French croisade.[185] Gibbon followed Thomas Fuller in dismissing the concept that the crusades were a legitimate defence, as they were disproportionate to the threat presented; Palestine was an objective, not because of reason but because of fanaticism and superstition.[186] William Robertson expanded on Fleury in a new, empirical, objective approach, placing crusading in a narrative of progress towards modernity. The cultural consequences of growth in trade, the rise of the Italian cities and progress are elaborated in his work. In this he influenced his student Walter Scott.[187] Much of the popular understanding of the crusades derives from the 19th century novels of Scott and the French histories by Joseph François Michaud.[188]
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+ In a 2001 article—"The Historiography of the Crusades"—Giles Constable attempted to categorise what is meant by "Crusade" into four areas of contemporary crusade study. His view was that Traditionalists such as Hans Eberhard Mayer are concerned with where the crusades were aimed, Pluralists such as Jonathan Riley-Smith concentrate on how the crusades were organised, Popularists including Paul Alphandery and Etienne Delaruelle focus on the popular groundswells of religious fervour, and Generalists, such as Ernst-Dieter Hehl focus on the phenomenon of Latin holy wars.[4][5] The historian Thomas F. Madden argues that modern tensions are the result of a constructed view of the crusades created by colonial powers in the 19th century and transmitted into Arab nationalism. For him the crusades are a medieval phenomenon in which the crusaders were engaged in a defensive war on behalf of their co-religionists.[189]
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+ The Muslim world exhibited little interest in the crusades until the middle of the 19th century. Arabic-speaking Syrian Christians began translating French histories into Arabic, leading to the replacement of the term "wars of the Ifranj" – Franks – with al-hurub al Salabiyya – "wars of the Cross". The Ottoman Turk Namık Kemal published the first modern Saladin biography in 1872. The Jerusalem visit in 1898 of Kaiser Wilhelm prompted further interest, with the Egyptian Sayyid Ali al-Hariri producing the first Arabic history of the crusades. Modern studies can be driven by political motives, such as the hope of learning from the Muslim forces' triumph over their enemies.[190]
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+ The Punic Wars were a series of three wars fought between Rome and Carthage from 264 BC to 146 BC. The main cause of the Punic Wars was the conflicts of interest between the existing Carthaginian Empire and the expanding Roman Republic. The Romans were initially interested in expansion via Sicily (which at that time was a cultural melting pot), part of which lay under Carthaginian control. At the start of the First Punic War (264–241 BC), Carthage was the dominant power of the Western Mediterranean, with an extensive maritime empire. Rome was a rapidly ascending power in Italy, but it lacked the naval power of Carthage. The Second Punic War (218–201 BC) witnessed Hannibal's crossing of the Alps in 218 BC, followed by a prolonged but ultimately failed campaign of Carthage's Hannibal in mainland Italy. By the end of the Third Punic War (149–146 BC), after more than a hundred years and the loss of many hundreds of thousands of soldiers from both sides, Rome had conquered Carthage's empire, completely destroyed the city, and became the most powerful state of the Western Mediterranean.
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+ With the end of the Macedonian Wars – which ran concurrently with the Punic Wars – and the defeat of the Seleucid King Antiochus III the Great in the Roman–Seleucid War (Treaty of Apamea, 188 BC) in the eastern sea, Rome emerged as the dominant Mediterranean power and one of the most powerful cities in classical antiquity. The Roman victories over Carthage in these wars gave Rome a preeminent status it would retain until the 5th century AD.
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+ The main source for almost every aspect of the Punic Wars[note 1] is the historian Polybius (c. 200 – c. 118 BC), a Greek sent to Rome in 167 BC as a hostage.[2] His works include a now-lost manual on military tactics,[3] but he is now known for The Histories, written sometime after 146 BC.[4][5] Polybius's work is considered broadly objective and largely neutral as between Carthaginian and Roman points of view.[6][7] Polybius was an analytical historian and wherever possible personally interviewed participants, from both sides, in the events he wrote about.[8][9][10] He accompanied the Roman general Scipio Aemilianus during his campaign in North Africa which resulted in a Roman victory in the Third Punic War.[11]
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+ The accuracy of Polybius's account has been much debated over the past 150 years, but the modern consensus is to accept it largely at face value, and the details of the war in modern sources are largely based on interpretations of Polybius's account.[2][12][13] The modern historian Andrew Curry sees Polybius as being "fairly reliable";[14] while Craige Champion describes him as "a remarkably well-informed, industrious, and insightful historian".[15]
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+ Other, later, ancient histories of the war exist, although often in fragmentary or summary form.[16] Modern historians usually take into account the writings of various Roman annalists, some contemporary; the Sicilian Greek Diodorus Siculus; the later Roman historians Livy (who relied heavily on Polybius[17]), Plutarch, Appian (whose account of the Third Punic War is especially valuable[18]) and Dio Cassius.[19] The classicist Adrian Goldsworthy states "Polybius' account is usually to be preferred when it differs with any of our other accounts".[note 2][9] Other sources include inscriptions, archaeological evidence and empirical evidence from reconstructions such as the trireme Olympias.[20]
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+ The Roman Republic had been aggressively expanding in the southern Italian mainland for a century before the First Punic War.[21] It had conquered peninsular Italy south of the River Arno by 272 BC, when the Greek cities of southern Italy (Magna Graecia) submitted at the conclusion of the Pyrrhic War.[22] During this period Carthage, with its capital in what is now Tunisia, had come to dominate southern Spain, much of the coastal regions of North Africa, the Balearic Islands, Corsica, Sardinia, and the western half of Sicily in a military and commercial empire.[23]
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+ Beginning in 480 BC, Carthage had fought a series of inconclusive wars against the Greek city states of Sicily, led by Syracuse.[24] By 264 BC Carthage and Rome were the preeminent powers in the western Mediterranean.[25] The two states had several times asserted their mutual friendship via formal alliances: in 509 BC, 348 BC and around 279 BC. Relationships were good, with strong commercial links. During the Pyrrhic War of 280–275 BC, against a king of Epirus who alternately fought Rome in Italy and Carthage on Sicily, Carthage provided materiel to the Romans and on at least one occasion used its navy to ferry a Roman force.[26][27]
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+ The two states had several times asserted their mutual friendship via formal alliances: in 509 BC, 348 BC and around 279 BC. Relationships were good, with strong commercial links. During the Pyrrhic War of 280–275 BC, against a king of Epirus who alternately fought Rome in Italy and Carthage on Sicily, Carthage provided materiel to the Romans and on at least one occasion used its navy to ferry a Roman force.[26][27] Rome's expansion into southern Italy probably made it inevitable that it would eventually clash with Carthage over Sicily on some pretext. The immediate cause of the war was the issue of control of the Sicilian town of Messana (modern Messina).[28] In 264 BC Carthage and Rome went to war, starting the First Punic War.[29]
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+ Most male Roman citizens were eligible for military service and would serve as infantry, a better-off minority providing a cavalry component. Traditionally, when at war the Romans would raise two legions, each of 4,200 infantry[note 3] and 300 cavalry. A few infantry served as javelin-armed skirmishers. The balance were equipped as heavy infantry, with body armour, a large shield and short thrusting swords. They were divided into three ranks, of which the front rank also carried two javelins, while the second and third ranks had a thrusting spear instead. Both legionary sub-units and individual legionaries fought in relatively open order. It was the long-standing Roman procedure to elect two men each year, known as consuls, to each lead an army. An army was usually formed by combining a Roman legion with a similarly sized and equipped legion provided by their Latin allies.[32][33]
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+ Carthaginian citizens only served in their army if there was a direct threat to the city. When they did they fought as well-armoured heavy infantry armed with long thrusting spears, although they were notoriously ill-trained and ill-disciplined. In most circumstances Carthage recruited foreigners to make up its army. Many would be from North Africa, a majority during the First Punic War, which provided several types of fighters including: close-order infantry equipped with large shields, helmets, short swords and long thrusting spears; javelin-armed light infantry skirmishers; close-order shock cavalry carrying spears; and light cavalry skirmishers who threw javelins from a distance and avoided close combat.[34][35][36] Both Spain and Gaul provided large numbers of experienced infantry, especially during the Second Punic War; they were mostly unarmoured troops who would charge ferociously, but had a reputation for breaking off if a combat was protracted.[37][38][39] The close order Libyan infantry and the citizen-militia would fight in a tightly packed formation known as a phalanx.[35] On occasion some of the infantry would wear captured Roman armour, especially among Hannibal's troops.[40] Slingers were frequently recruited from the Balearic Islands.[41][42] The Carthaginians also employed war elephants; North Africa had indigenous African forest elephants at the time.[note 4][38][44]
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+ Garrison duty and land blockades were the most common operations for both armies.[45][46] When armies were campaigning, surprise attacks, ambushes and strategems were common.[35][47] More formal battles were usually preceded by the two armies camping a mile or two apart (2–3 km) for days or weeks; sometimes forming up in battle order each day. If neither commander could see an advantage, both sides might march off without engaging.[48] Forming up in battle order was a complicated and premeditated affair, which took several hours. Infantry were usually positioned in the centre of the battle line, with light infantry skirmishers to their front and cavalry on each flank.[49] Many battles were decided when one side's infantry force was attacked in the flank or rear and they were partially or wholly enveloped.[35]
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+ Quinqueremes, meaning "five-oarsmen",[50] provided the workhorses of the Roman and Carthaginian fleets throughout the Punic Wars.[51] So ubiquitous was the type that Polybius uses it as a shorthand for "warship" in general.[52] A quinquereme carried a crew of 300: 280 oarsmen and 20 deck crew and officers.[53] It would also normally carry a complement of 40 marines,[54] if battle was thought to be imminent this would be increased to as many as 120.[55][56] In 260 BC Romans set out to construct a fleet and used a shipwrecked Carthaginian quinquereme as a blueprint for their own.[57]
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+ As novice shipwrights, the Romans built copies that were heavier than the Carthaginian vessels, and so slower and less manoeuvrable.[58] Getting the oarsmen to row as a unit, let alone to execute more complex battle manoeuvres, required long and arduous training.[59] At least half of the oarsmen would need to have had some experience if the ship was to be handled effectively.[60] As a result, the Romans were initially at a disadvantage against the more experienced Carthaginians. To counter this, the Romans introduced the corvus, a bridge 1.2 metres (4 feet) wide and 11 metres (36 feet) long, with a heavy spike on the underside, which was designed to pierce and anchor into an enemy ship's deck.[55] This allowed Roman legionaries acting as marines to board enemy ships and capture them, rather than employing the previously traditional tactic of ramming.[61]
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+ All warships were equipped with rams, a triple set of 60-centimetre-wide (2 ft) bronze blades weighing up to 270 kilograms (600 lb) positioned at the waterline. In the century prior to the Punic Wars, boarding had become increasingly common and ramming had declined, as the larger and heavier vessels adopted in this period lacked the speed and manoeuvrability necessary to ram, while their sturdier construction reduced the ram's effect even in case of a successful attack. The Roman adaptation of the corvus was a continuation of this trend and compensated for their initial disadvantage in ship-manoeuvring skills. The added weight in the prow compromised both the ship's manoeuvrability and its seaworthiness, and in rough sea conditions the corvus became useless; part way through the First Punic War the Romans ceased using it.[61][62][63]
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+ The war began with the Romans gaining a foothold on Sicily at Messana (modern Messina).[64] The Romans then pressed Syracuse, the only significant independent power on the island, into allying with them[65] and laid siege to Carthage's main base at Akragas.[66] A large Carthaginian army attempted to lift the siege in 262 BC, but was heavily defeated at the Battle of Akragas. That night the Carthaginian garrison escaped and the Romans seized the city and its inhabitants, selling 25,000 of them into slavery.[67] The Romans then built a navy to challenge Carthage's,[68] and using the corvus inflicted several defeats.[69][70][71] A Carthaginian base on Corsica was seized, but an attack on Sardinia was repulsed; the base on Corsica was then lost.[72]
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+ Taking advantage of their naval victories the Romans launched an invasion of North Africa,[73] which the Carthaginians intercepted. At the Battle of Cape Ecnomus the Carthaginians were again beaten;[74] this was possibly the largest naval battle in history by the number of combatants involved.[75][76][77] The invasion initially went well and in 255 BC the Carthaginians sued for peace; the proposed terms were so harsh they fought on,[78] defeating the invaders.[79] The Romans sent a fleet to evacuate their survivors and the Carthaginians opposed it at the Battle of Cape Hermaeum off Africa; the Carthaginians were heavily defeated.[80] The Roman fleet, in turn, was devastated by a storm while returning to Italy, losing most of its ships and over 100,000 men.[80][81][82]
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+ The war continued, with neither side able to gain a decisive advantage.[83] The Carthaginians attacked and recaptured Akragas in 255 BC, but not believing they could hold the city, they razed and abandoned it.[84][85] The Romans rapidly rebuilt their fleet, adding 220 new ships, and captured Panormus (modern Palermo) in 254 BC.[86] The next year they lost another 150 ships to a storm.[87] In 251 BC the Carthaginians attempted to recapture Panormus, but were defeated in a battle outside the walls.[88][89] Slowly the Romans had occupied most of Sicily; in 249 BC they besieged the last two Carthaginian strongholds – in the extreme west.[90] They also launched a surprise attack on the Carthaginian fleet, but were defeated at the Battle of Drepana.[91] The Carthaginians followed up their victory and most of the remaining Roman warships were lost at the Battle of Phintias.[92] After several years of stalemate,[93] the Romans rebuilt their fleet again in 243 BC[94] and effectively blockaded the Carthaginian garrisons[95]. Carthage assembled a fleet which attempted to relieve them, but it was destroyed at the Battle of the Aegates Islands in 241 BC,[96][97] forcing the cut-off Carthaginian troops on Sicily to negotiate for peace.[98][95]
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+ A treaty was agreed. By its terms Carthage paid 3,200 talents of silver[note 5][note 6] in reparations and Sicily was annexed as a Roman province.[96] Henceforth Rome considered itself the leading military power in the western Mediterranean, and increasingly the Mediterranean region as a whole. The immense effort of repeatedly building large fleets of galleys during the war laid the foundation for Rome's maritime dominance for 600 years.[100]
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+ The Mercenary War began in 241 BC as a dispute over the payment of wages owed to 20,000 foreign soldiers who had fought for Carthage on Sicily during the First Punic War. When a compromise seemed to have been reached the army erupted into full-scale mutiny under the leadership of Spendius and Matho. 70,000 Africans from Carthage's oppressed dependant territories flocked to join them, bringing supplies and finance. War-weary Carthage fared poorly in the initial engagements, especially under the generalship of Hanno. Hamilcar Barca, a veteran of the campaigns in Sicily (and father of Hannibal Barca), was given joint command of the army in 240 BC; and supreme command in 239 BC. He campaigned successfully, initially demonstrating leniency in an attempt to woo the rebels over. To prevent this, in 240 BC Spendius tortured 700 Carthaginian prisoners to death, and henceforth the war was pursued with great brutality.
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+ By early 237 BC, after numerous setbacks, the rebels were defeated and their cities brought back under Carthaginian rule. An expedition was prepared to reoccupy Sardinia, where mutinous soldiers had slaughtered all Carthaginians. The Roman Senate stated they considered the preparation of this force an act of war, and demanded Carthage cede Sardinia and Corsica, and pay an additional 1,200-talent indemnity.[101][102][note 7] Weakened by 30 years of war, Carthage agreed rather than again enter into conflict with Rome.[103] Polybius considered this "contrary to all justice"[101] and modern historians have variously described the Romans' behaviour as "unprovoked aggression and treaty-breaking",[101] "shamelessly opportunistic"[104] and an "unscrupulous act".[105] These events fuelled resentment in Carthage, which was not reconciled to Rome's perception of its situation. This breach of the recently signed treaty has been considered to be the single greatest cause of war with Carthage breaking out again in 218 BC in the Second Punic War.[106][107][108]
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+
45
+ With the suppression of the rebellion, Hamilcar Barca understood that Carthage needed to strengthen its economic and military base if it were to again confront Rome.[110] After the First Punic War, Carthaginian possessions in Iberia (modern Spain and Portugal) were limited to a handful of prosperous coastal cities.[111] Hamilcar took the army which he had led to victory in the Mercenary War and led it to Iberia in 237 BC. He carved out a quasi-monarchial, autonomous state in south-east Iberia.[112] This gave Carthage the silver mines, agricultural wealth, manpower, military facilities such as shipyards and territorial depth to stand up to future Roman demands with confidence.[113][114] Hamilcar ruled as a viceroy and was succeeded by his son-in-law, Hasdrubal, in the early 220s BC and then his son, Hannibal, in 221 BC.[115] In 226 BC the Ebro Treaty was agreed, specifying the Ebro River as the northern boundary of the Carthaginian sphere of influence.[116] A little later Rome made a separate treaty with the city of Saguntum, well south of the Ebro.[117]
46
+
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+ In 219 BC a Carthaginian army under Hannibal besieged, captured and sacked Saguntum.[106][118] In spring 218 BC Rome declared war on Carthage.[119] There were three main military theatres in the war: Italy, where Hannibal defeated the Roman legions repeatedly, with occasional subsidiary campaigns in Sicily and Greece; Iberia, where Hasdrubal, a younger brother of Hannibal, defended the Carthaginian colonial cities with mixed success until moving into Italy; and Africa, where the war was decided.
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+ In 218 BC there was some naval skirmishing in the waters around Sicily. The Romans beat off a Carthaginian attack[120][121] and captured the island of Malta.[122] In Cisalpine Gaul (modern northern Italy), the major Gallic tribes attacked the Roman colonies, causing the Romans to flee to Mutina (modern Modena), which the they besieged. A Roman relief army raised the siege, but was then ambushed and besieged itself.[123] An army had been raised to campaign in Iberia under the brothers Gnaeus and Publius Scipio and the Roman Senate detached one Roman and one allied legion from it tosend to the region. The Scipios had to raise fresh troops to replace these and thus could not set out for Iberia until September.[124]
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+
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+ Meanwhile, Hannibal assembled a Carthaginian army in New Carthage (modern Cartagena) and led it northwards along the coast in May or June. It entered Gaul and took an inland route, to avoid the Roman allies along the coast.[125] At the Battle of Rhone Crossing, Hannibal defeated a force of local Allobroges that sought to bar his way.[126] A Roman fleet carrying the Scipio brothers' army landed at Rome's ally Massalia (modern Marseille) at the mouth of the Rhone.[127] Hannibal evaded the Romans and Gnaeus Scipio continued to Iberia with the Roman army;[128][129] Publius returned to Rome.[129] The Carthaginians reached the foot of the Alps by late autumn[125] and crossed them, surmounting the difficulties of climate, terrain[125]> and the guerrilla tactics of the native tribes. The exact route is disputed. Hannibal arrived with 20,000 infantry, 6,000 cavalry, and an unknown number of elephants[65] in what is now Piedmont, northern Italy. The Romans were still in their winter quarters. His surprise entry into the Italian peninsula led to the termination of Rome's planned campaign for the year, an invasion of Africa.[130]
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+
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+ Hannibal's first action was to take the chief city of the hostile Taurini (in the area of modern-day Turin). His army then routed the cavalry and light infantry of the Romans under Publius Scipio at the Battle of Ticinus.[131] As a result, most of the Gallic tribes declared for the Carthaginian cause, and Hannibal's army grew to over 40,000 men.[132] A large Roman army under the command of Sempronius Longus was lured into combat by Hannibal at the Battle of the Trebia, encircled and destroyed.[133] Only 10,000 Romans out of 42,000 were able to cut their way to safety. Gauls now joined Hannibal's army in large numbers, bringing it up to 60,000 men.[132] The Romans stationed an army at Arretium and one on the Adriatic coast to block Hannibal's advance into central Italy.[134]
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+
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+ In early spring 217 BC, the Carthaginians crossed the Apennines unopposed, taking a difficult but unguarded route.[135] Hannibal attempted without success to draw the main Roman army under Gaius Flaminius into a pitched battle by devastating the area they had been sent to protect.[136] Hannibal then cut off the Roman army from Rome, which provoked Flaminius into a hasty pursuit without proper reconnaissance.[137] Then, in a defile on the shore of Lake Trasimenus, Hannibal set an ambush[137] and in the Battle of Lake Trasimene completely defeated the Roman army and killed Flaminius.[137] 15,000 Romans were killed and 15,000 taken prisoner. 4,000 Roman cavalry from their other army were also engaged and wiped out.[138] The prisoners were sold as slaves if they were Romans, but released if they were from one of Rome's Latin allies.[139] Hannibal hoped that some of these allies could be persuaded to defect, and marched south in the hope of winning over allies among the ethnic Greek and Italic city states.[134]
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+
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+ The Romans, panicked by these heavy defeats, appointed Quintus Fabius Maximus as dictator.[139] Fabius invented the Fabian strategy of avoiding open battle with his opponent, but constantly skirmishing with small detachments of the enemy. This was not popular among the soldiers, the Roman public nor the Roman elite, since he avoided battle while Italy was being devastated by the enemy.[134] Hannibal marched through the richest and most fertile provinces of Italy, hoping that the devastation would draw Fabius into battle, but Fabius refused.[140]
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+ At the elections of 216 BC the more aggressive minded Gaius Terentius Varro and Lucius Aemilius Paullus were elected as consuls.[141] The Roman Senate authorized the raising of double-sized armies, a force of 86,000 men, the largest in Roman history up to that point.[141] Paullus and Varro marched southward to confront Hannibal, who accepted battle on the open plain near Cannae. In the Battle of Cannae The Roman legions forced their way through Hannibal's deliberately weak centre, but the Libyans on the wings swung around their advance, menacing their flanks.[142] Hasdrubal led Carthaginian cavalry on the left wing and routed the Roman cavalry opposite, then swept around the rear of the Romans to attack the cavalry on the other wing. He then charged into the legions' from behind.[142] As a result, the Roman infantry was surrounded with no means of escape.[142] At least 67,500 Romans were killed or captured.[142].
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+ Gnaeus Scipio continued on from Massala in the summer of 218 BC to Iberia (modern Spain and Portugal), where he won support among the local tribes.[128] The Carthaginian commander in the area refused to wait for reinforcements and attacked Scipio at the Battle of Cissa in late 218 BC and was defeated.[128][143] In 217 BC, the Carthaginians moved to engage the combined Roman and Massalian fleet at the Battle of Ebro River. The 40 Carthaginian and Iberian vessels were beaten by 55 Roman and Massalian ships in the second naval engagement of the war, with 29 Carthaginian ships lost. Carthaginian forces retreated, but the Romans remained confined to the area between the Ebro and Pyrenees.[143]
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+
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+ The Roman army in Spain was preventing the Carthaginians from sending reinforcements from Iberia to Hannibal or to the insurgent Gauls in northern Italy.[143] Hasdrubal marched into Roman territory in 215 BC, besieged a pro-Roman town and offered battle at Dertosa. After a hard-fought battle, he was defeated although both sides suffered heavy losses.[144] Hasdrubal was now unable to reinforce Hannibal in Italy.[144][128]
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+ The Carthaginians suffered a wave of defections of local Celtiberian tribes to Rome.[128] The Scipio brothers captured Saguntum in 212 BC.[144] In 211, they hired 20,000 Celtiberian mercenaries to reinforce their army.[144] Observing that the three Carthaginian armies were deployed apart from each other, the Scipios split their forces.[144] Publius moved to attack Mago Barca near Castulo, while Gnaeus marched on Hasdrubal.[144] This stratagy resulted in the Battle of Castulo and the Battle of Ilorca, usually combined as the Battle of the Upper Baetis.[128][144] Both battles ended in complete defeat for the Romans, with both of the Scipio brothers being killed, as Hasdrubal had bribed the Romans' mercenaries to desert.[128][144] The Romans retreated to their coastal stronghold north of the Ebro, from which the Carthaginians failed to expel them.[144][128] Claudius Nero brought over reinforcements in 210 and stabilized the situation.[144]
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+ In 210 BC, Scipio Africanus arrived in Spain with further reinforcements.[145] In a carefully planned assault in 209 BC, he captured the lightly-defended centre of Carthaginian power in Spain, Cartago Nova.[146][145] Scipio had the population slaughtered and a vast booty of gold, silver and siege artillery was taken.[147][145] He liberated the Iberian hostages kept by the Carthaginians to ensure the loyalty of the Iberian tribes,[147][145] although many of them were subsequently to fight against the Romans.[145]
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+
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+ In 206 BC, at the Battle of Ilipa, Scipio with 48,000 men, half Italians and half Iberians, defeated a Carthaginian army of 54,500 men and 32 elephants under the command of Mago Barca, Hasdrubal Gisco and Masinissa This sealed the fate of the Carthaginian presence in Iberia.[148][145] It was followed by the Roman capture of Gades after the city rebelled against Carthaginian rule.[149]
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+ Later that year a dangerous mutiny broke out among Roman troops at their camp at Sucro. It initially attracted support from Iberian leaders, disappointed that Roman forces had remained in the peninsula after the expulsion of the Carthaginians. It was effectively put down by Scipio Africanus. In 205 BC a last attempt was made by Mago to recapture New Carthage when the Roman occupiers were shaken by another mutiny and an Iberian uprising, but he was repulsed. Mago left Spain for Italy with his remaining forces.[147] In 203 BC Carthage succeeded in recruiting at least 4,000 mercenaries from Iberia, despite Rome's nominal control.[150]
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+ In 213 BC Syphax, a powerful Numidian king in North Africa,[144] declared for Rome. Rome sent advisers to train his soldiers[144] and he waged war against the Carthaginian ally Gala.[144] In 206 BC the Carthaginians ended this drain on their resources by dividing several Numidian kingdoms with him. One of those disinherited was the Numidian prince Masinissa, who was thus driven into the arms of Rome.[151]
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+ In 205 BC Scipio Africanus was given command of the legions in Sicily and allowed to enroll volunteers for his plan to end the war by an invasion of Africa.[152] After landing in Africa in 204 BC, he was joined by Masinissa and a force of Numidian cavalry.[153] Scipio then besieged and failed to take the city of Utica.[154] When a Carthaginian and Numidian relief army under Hasdrubal Barca and Syphax moved to confront him, he mounted a surprise attack and destroyed it.[155] In 203 BC Scipio confronted a second Carthaginian army and destroyed at the Battle of the Great Plains. King Syphax was pursued and taken prisoner at the Battle of Cirta and Masinissa seized a large part of his kingdom with Roman help.[156]
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+ Rome and Carthage entered into peace negotiations, and Carthage recalled Hannibal from Italy.[157] Largely due to mutual mistrust the negotiations came to nothing.[158] Hannibal was placed in command of another army, based his veterans from Italy and newly raised troops from Africa, but with few cavalry.[159] The decisive Battle of Zama followed in October 202 BC.[160] Unlike most battles of the Second Punic War, the Romans had superiority in cavalry and the Carthaginians in infantry.[159] Hannibal attempted to use 80 elephants to break into the Roman infantry formation, but the Romans countered them effectively and they routed back through the Carthaginian ranks.[161] The Roman and Allied Numidian cavalry drove the Carthaginian cavalry from the field. The two sides' infantry fought inconclusively until the Roman cavalry returned and attacked his rear. The Carthaginian formation collapsed; Hannibal was one of the few to escape the field.[160]
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+ The peace treaty imposed on the Carthaginians stripped them of all of their overseas territories, and some of their African ones. An indemnity of 10,000 silver talents[note 8] was to be paid over 50 years. Hostages were taken. Carthage was forbidden to possess war elephants and its fleet was restricted to 10 warships. It was prohibited from waging war outside Africa, and in Africa only with Rome's express permission. Many senior Carthaginians wanted to reject it, but Hannibal spoke strongly in its favour and it was accepted in spring 201 BC.[162] Henceforth it was clear that Carthage was politically subordinate to Rome.[163]
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+ At the end of the war, Masinissa emerged as by far the most powerful ruler among the Numidians.[164] Over the following 48 years he repeatedly took advantage of Carthage's inability to protect its possessions. Whenever Carthage petioned Rome for redress, or permission to take military action, Rome backed its ally, Masinissa, and refused.[165] Masinissa's seizures of and raids into Carthaginian territory became increasingly flagrant. In 151 BC Carthage raised a large army, the treaty notwithstanding, and counter attacked the Numidians. The campaign ended in disaster and the army surrendered.[166] Carthage had paid off its indemnity and was prospering economically, but was no military threat to Rome.[167][168] Elements in the Roman Senate had long wished to destroy Carthage, and with the breach of the treaty as a casus belli, war was declared in 149 BC.[166]
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+ In 149 BC a Roman army of approximately 50,000 men, jointly commanded by both consuls, landed near Utica, 35 kilometres (22 mi) north of Carthage.[169] Rome demanded that if war were to be avoided, the Carthaginians must hand over all of their armaments. Vast amounts of materiel delivered, including 200,000 sets of armour, 2,000 catapults and a large number of warships.[170] This done, the Romans demanded that the Carthaginians burn their city and relocate at least 10 miles (16 km) from the sea; the Carthaginians broke off negotiations and set to recreating their armoury.[171]
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+ As well as manning the walls of Carthage, the Carthaginians formed a field army under Hasdrubal, which was based 25 kilometres (16 mi) to the south.[173][174] The Roman army moved to lay siege to Carthage, but its walls were so strong and its citizen-militia so determined that it was unable to make any impact, while the Carthaginians struck back effectively. Their army effectively raided the Roman lines of communication,[174] and in 148 BC Carthaginian fire ships destroyed many Roman vessels. The main Roman camp was in a swamp, which caused an outbreak of disease during the summer.[175] The Romans moved their camp, and their ships, further away; so that they were more blockading than closely besieging the city.[176] The war dragged on into 147 BC.[174]
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+ In early 147 BC Scipio Aemilianus, an adopted grandson of Scipio Africanus who had distinguished himself during the previous two years' fighting, was elected consul and took control of the war.[166][177] The Carthaginians continued to resist vigorously: they constructed warships and during the summer twice gave battle to the Roman fleet, losing both times.[177] The Romans launched an assault on the walls; after confused fighting they broke into the city, but lost in the dark, withdrew. Hasdrubal and his army withdrew into the city to reinforce the garrison.[178] Roman prisoners were taken and Hasdrubal had them tortured to death on the walls, in view of the Roman army. He was reinforcing the will to resist in the Carthaginian citizens; from this point there could be no possibility of negotiations. Some members of the city council denounced his actions and Hasdrubal had them too put to death and took over control of the city.[177][179] With no Carthaginian army in the field those cities which had remained loyal went over to the Romans or were captured.[180]
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+ Scipio moved back to a close blockade of the city, and built a mole which cut off supply from the sea.[181] In the spring of 146 BC the Roman army managed to secure a foothold on the fortifications near the harbour.[182][183] When the main assault began it quickly captured the cities main square, where the legions camped overnight.[184] The next morning the Romans systematically worked their way through the residential part of the city, killing everyone they encountered and firing the buildings behind them.[182] At times the Romans progressed from rooftop to rooftop, to prevent missiles being hurled down on them.[184] It took six days to clear the city of resistance, and on the last day Scipio agreed to accept prisoners. The last holdouts, including Roman deserters in Carthaginian service, fought on from the Temple of Eshmoun and burnt it down around themselves when all hope was gone.[185] The 50,000 Carthaginians, a small part of the pre-war population, were sold into slavery.[186] The notion that Roman forces then sowed the city with salt to ensure that nothing would grow there again is a 20th-century invention.[187]
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+ The remaining Carthaginian territories were annexed by Rome and reconstituted to become the Roman province of Africa. Numerous significant Punic cities, such as those in Mauretania, were taken over and rebuilt by the Romans.[188] Utica, the Punic city which changed loyalties at the beginning of the siege, became the capital of the Roman province of Africa.[189] A century later, the site of Carthage was rebuilt as a Roman city by Julius Caesar, and would become one of the main cities of Roman Africa by the time of the Empire.
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+ Rome still exists as the capital of Italy; the ruins of Carthage lie 16 kilometres (10 mi) east of modern Tunis on the North African coast.[190]
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+ An airport is an aerodrome with extended facilities, mostly for commercial air transport.[1][2] Airports often have facilities to store and maintain aircraft, and a control tower. An airport consists of a landing area, which comprises an aerially accessible open space including at least one operationally active surface such as a runway for a plane to take off[3] or a helipad,[4] and often includes adjacent utility buildings such as control towers, hangars[5] and terminals. Larger airports may have airport aprons, taxiway bridges, air traffic control centres, passenger facilities such as restaurants and lounges, and emergency services. In some countries, the US in particular, airports also typically have one or more fixed-base operators, serving general aviation.
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+ An airport solely serving helicopters is called a heliport. An airport for use by seaplanes and amphibious aircraft is called a seaplane base. Such a base typically includes a stretch of open water for takeoffs and landings, and seaplane docks for tying-up.
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+ An international airport has additional facilities for customs and passport control as well as incorporating all the aforementioned elements. Such airports rank among the most complex and largest of all built typologies, with 15 of the top 50 buildings by floor area being airport terminals.[citation needed][6]
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+ The terms aerodrome, airfield, and airstrip also refer to airports, and the terms heliport, seaplane base, and STOLport refer to airports dedicated exclusively to helicopters, seaplanes, and short take-off and landing aircraft.
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+ In colloquial use in certain environments, the terms airport and aerodrome are often interchanged. However, in general, the term airport may imply or confer a certain stature upon the aviation facility that other aerodromes may not have achieved. In some jurisdictions, airport is a legal term of art reserved exclusively for those aerodromes certified or licensed as airports by the relevant national aviation authority after meeting specified certification criteria or regulatory requirements.[7]
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+ That is to say, all airports are aerodromes, but not all aerodromes are airports. In jurisdictions where there is no legal distinction between aerodrome and airport, which term to use in the name of an aerodrome may be a commercial decision. In US technical/legal usage, landing area is used instead of aerodrome, and airport means "a landing area used regularly by aircraft for receiving or discharging passengers or cargo".[8]
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+ Smaller or less-developed airfields, which represent the vast majority, often have a single runway shorter than 1,000 m (3,300 ft). Larger airports for airline flights generally have paved runways of 2,000 m (6,600 ft) or longer. Skyline Airport in Inkom, Idaho has a runway that is only 122 m (400 ft) long.[9]
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+ In the United States, the minimum dimensions for dry, hard landing fields are defined by the FAR Landing And Takeoff Field Lengths. These include considerations for safety margins during landing and takeoff.
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+ The longest public-use runway in the world is at Qamdo Bamda Airport in China. It has a length of 5,500 m (18,045 ft). The world's widest paved runway is at Ulyanovsk Vostochny Airport in Russia and is 105 m (344 ft) wide.
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+ As of 2009[update], the CIA stated that there were approximately 44,000 "airports or airfields recognizable from the air" around the world, including 15,095 in the US, the US having the most in the world.[10][11]
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+ Most of the world's large airports are owned by local, regional, or national government bodies who then lease the airport to private corporations who oversee the airport's operation. For example, in the UK the state-owned British Airports Authority originally operated eight of the nation's major commercial airports – it was subsequently privatized in the late 1980s, and following its takeover by the Spanish Ferrovial consortium in 2006, has been further divested and downsized to operating just Heathrow. Germany's Frankfurt Airport is managed by the quasi-private firm Fraport. While in India GMR Group operates, through joint ventures, Indira Gandhi International Airport and Rajiv Gandhi International Airport. Bengaluru International Airport and Chhatrapati Shivaji International Airport are controlled by GVK Group. The rest of India's airports are managed by the Airports Authority of India. In Pakistan nearly all civilian airports are owned and operated by the Pakistan Civil Aviation Authority except for Sialkot International Airport which has the distinction of being the first privately owned public airport in Pakistan and South Asia[citation needed].
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+ In the US, commercial airports are generally operated directly by government entities or government-created airport authorities (also known as port authorities), such as the Los Angeles World Airports authority that oversees several airports in the Greater Los Angeles area, including Los Angeles International Airport[citation needed].
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+ In Canada, the federal authority, Transport Canada, divested itself of all but the remotest airports in 1999/2000. Now most airports in Canada are owned and operated by individual legal authorities or are municipally owned.
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+ Many US airports still lease part or all of their facilities to outside firms, who operate functions such as retail management and parking. All US commercial airport runways are certified by the FAA[12] under the Code of Federal Regulations Title 14 Part 139, "Certification of Commercial Service Airports"[13] but maintained by the local airport under the regulatory authority of the FAA.
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+ Despite the reluctance to privatize airports in the US (contrary to the FAA sponsoring a privatization program since 1996), the government-owned, contractor-operated (GOCO) arrangement is the standard for the operation of commercial airports in the rest of the world.
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+ The Airport & Airway Trust Fund (AATF) was created by the Airport and Airway Development in 1970 which finances aviation programs in the United States.[14] Airport Improvement Program (AIP), Facilities and Equipment (F&E), and Research, Engineering, and Development (RE&D) are the three major accounts of Federal Aviation Administration which are financed by the AATF, as well as pays for the FAA's Operation and Maintenance (O&M) account.[15] The funding of these accounts are dependent on the taxes the airports generate of revenues. Passenger tickets, fuel, and cargo tax are the taxes that are paid by the passengers and airlines help fund these accounts.[16]
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+ Airports revenues are divided into three major parts: aeronautical revenue, non-aeronautical revenue, and non-operating revenue. Aeronautical revenue makes up 56%, non-aeronautical revenue makes up 40%, and non-operating revenue makes up 4% of the total revenue of airports.[17]
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+ Aeronautical revenue are generated through airline rents and landing, passenger service, parking, and hangar fees. Landing fees are charged per aircraft for landing an airplane in the airport property.[18] Landing fees are calculated through the landing weight and the size of the aircraft which varies but most of the airports have a fixed rate and a charge extra for extra weight.[19] Passenger service fees are charges per passengers for the facilities used on a flight like water, food, wifi and shows which is paid while paying for an airline ticket.[citation needed] Aircraft parking is also a major revenue source for airports. Aircraft are parked for a certain amount of time before or after takeoff and have to pay to park there.[20] Every airport has his own rates of parking but at John F Kennedy airport in New York City charges $45 per hour for the plane of 100,000 pounds and the price increases with weight.[21]
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+ Non-aeronautical revenue is gained through things other than aircraft operations. It includes lease revenue from compatible land-use development, non-aeronautical building leases, retail and concession sales, rental car operations, parking and in-airport advertising.[22] Concession revenue is one big part of non-aeronautical revenue airports makes through duty free, bookstores, restaurants and money exchange.[20] Car parking is a growing source of revenue for airports, as more people use the parking facilities of the airport. O'Hare International Airport in Chicago charges $2 per hour for every car.[23]
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+ Airports are divided into landside and airside areas. The landside area is open to the public, while access to the airside area is tightly controlled. The airside area includes all parts of the airport around the aircraft, and the parts of the buildings that are accessible only to passengers and staff. Passengers and staff must be checked by security before being permitted to enter the airside area. Conversely, passengers arriving from an international flight must pass through border control and customs to access the landside area, where they can exit the airport. Many major airports will issue a secure keycard called an airside pass to employees, as some roles require employees to frequently move back and forth between landside and airside as part of their duties.
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+ A terminal is a building with passenger facilities. Small airports have one terminal. Large ones often have multiple terminals, though some large airports like Amsterdam Airport Schiphol still have one terminal. The terminal has a series of gates, which provide passengers with access to the plane.
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+ The following facilities are essential for departing passengers:
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+ The following facilities are essential for arriving passengers:
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+ For both sets of passengers, there must be a link between the passenger facilities and the aircraft, such as jet bridges or airstairs. There also needs to be a baggage handling system, to transport baggage from the baggage drop-off to departing planes, and from arriving planes to the baggage reclaim.
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+ The area where the aircraft park to load passengers and baggage is known as an apron or ramp (or incorrectly[24], "the tarmac").
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+ Airports with international flights have customs and immigration facilities. However, as some countries have agreements that allow travel between them without customs and immigrations, such facilities are not a definitive need for an international airport. International flights often require a higher level of physical security, although in recent years, many countries have adopted the same level of security for international and domestic travel.
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+ "Floating airports" are being designed which could be located out at sea and which would use designs such as pneumatic stabilized platform technology.
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+ Airport security normally requires baggage checks, metal screenings of individual persons, and rules against any object that could be used as a weapon. Since the September 11 attacks and the Real ID Act of 2005, airport security has dramatically increased and got tighter and stricter than ever before.
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+ Most major airports provide commercial outlets for products and services. Most of these companies, many of which are internationally known brands, are located within the departure areas. These include clothing boutiques and restaurants and in the US amounted to $4.2 billion in 2015.[25] Prices charged for items sold at these outlets are generally higher than those outside the airport. However, some airports now regulate costs to keep them comparable to "street prices". This term is misleading as prices often match the manufacturers' suggested retail price (MSRP) but are almost never discounted.[citation needed]
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+ Apart from major fast food chains, some airport restaurants offer regional cuisine specialties for those in transit so that they may sample local food or culture without leaving the airport.[26]
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+ Some airport structures include on-site hotels built within or attached to a terminal building. Airport hotels have grown popular due to their convenience for transient passengers and easy accessibility to the airport terminal. Many airport hotels also have agreements with airlines to provide overnight lodging for displaced passengers.
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+ Major airports in such countries as Russia and Japan offer miniature sleeping units within the airport that are available for rent by the hour. The smallest type is the capsule hotel popular in Japan. A slightly larger variety is known as a sleep box. An even larger type is provided by the company YOTEL.
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+ Airports may also contain premium and VIP services. The premium and VIP services may include express check-in and dedicated check-in counters.
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+ These services are usually reserved for first and business class passengers, premium frequent flyers, and members of the airline's clubs. Premium services may sometimes be open to passengers who are members of a different airline's frequent flyer program. This can sometimes be part of a reciprocal deal, as when multiple airlines are part of the same alliance, or as a ploy to attract premium customers away from rival airlines.
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+ Sometimes these premium services will be offered to a non-premium passenger if the airline has made a mistake in handling of the passenger, such as unreasonable delays or mishandling of checked baggage.
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+ Airline lounges frequently offer free or reduced cost food, as well as alcoholic and non-alcoholic beverages. Lounges themselves typically have seating, showers, quiet areas, televisions, computer, Wi-Fi and Internet access, and power outlets that passengers may use for their electronic equipment. Some airline lounges employ baristas, bartenders and gourmet chefs.
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+ Airlines sometimes operate multiple lounges within the one airport terminal allowing ultra-premium customers, such as first class customers, additional services, which are not available to other premium customers. Multiple lounges may also prevent overcrowding of the lounge facilities.
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+ In addition to people, airports move cargo around the clock. Cargo airlines often have their own on-site and adjacent infrastructure to transfer parcels between ground and air.
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+ Cargo Terminal Facilities are areas where international airports export cargo has to be stored after customs clearance and prior to loading on the aircraft. Similarly import cargo that is offloaded needs to be in bond before the consignee decides to take delivery. Areas have to be kept aside for examination of export and import cargo by the airport authorities. Designated areas or sheds may be given to airlines or freight forward ring agencies.
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+ Every cargo terminal has a landside and an airside. The landside is where the exporters and importers through either their agents or by themselves deliver or collect shipments while the airside is where loads are moved to or from the aircraft. In addition cargo terminals are divided into distinct areas – export, import and interline or transshipment.
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+ Airports require parking lots, for passengers who may leave the cars at the airport for a long period of time. Large airports will also have car-rental firms, taxi ranks, bus stops and sometimes a train station.
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+
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+ Many large airports are located near railway trunk routes for seamless connection of multimodal transport, for instance Frankfurt Airport, Amsterdam Airport Schiphol, London Heathrow Airport, Tokyo Haneda Airport, Tokyo Narita Airport, London Gatwick Airport and London Stansted Airport. It is also common to connect an airport and a city with rapid transit, light rail lines or other non-road public transport systems. Some examples of this would include the AirTrain JFK at John F. Kennedy International Airport in New York, Link Light Rail that runs from the heart of downtown Seattle to Seattle–Tacoma International Airport, and the Silver Line T at Boston's Logan International Airport by the Massachusetts Bay Transportation Authority (MBTA). Such a connection lowers risk of missed flights due to traffic congestion. Large airports usually have access also through controlled-access highways ('freeways' or 'motorways') from which motor vehicles enter either the departure loop or the arrival loop.
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+ The distances passengers need to move within a large airport can be substantial. It is common for airports to provide moving walkways, buses, and rail transport systems. Some airports like Hartsfield–Jackson Atlanta International Airport and London Stansted Airport have a transit system that connects some of the gates to a main terminal. Airports with more than one terminal have a transit system to connect the terminals together, such as John F. Kennedy International Airport, Mexico City International Airport and London Gatwick Airport.
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+ There are three types of surface that aircraft operate on:
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+ Air traffic control (ATC) is the task of managing aircraft movements and making sure they are safe, orderly and expeditious. At the largest airports, air traffic control is a series of highly complex operations that requires managing frequent traffic that moves in all three dimensions.
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+ A "towered" or "controlled" airport has a control tower where the air traffic controllers are based. Pilots are required to maintain two-way radio communication with the controllers, and to acknowledge and comply with their instructions. A "non-towered" airport has no operating control tower and therefore two-way radio communications are not required, though it is good operating practice for pilots to transmit their intentions on the airport's common traffic advisory frequency (CTAF) for the benefit of other aircraft in the area. The CTAF may be a Universal Integrated Community (UNICOM), MULTICOM, Flight Service Station (FSS), or tower frequency.
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+ The majority of the world's airports are small facilities without a tower. Not all towered airports have 24/7 ATC operations. In those cases, non-towered procedures apply when the tower is not in use, such as at night. Non-towered airports come under area (en-route) control. Remote and virtual tower (RVT) is a system in which ATC is handled by controllers who are not present at the airport itself.
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+ Air traffic control responsibilities at airports are usually divided into at least two main areas: ground and tower, though a single controller may work both stations. The busiest airports may subdivide responsibilities further, with clearance delivery, apron control, and/or other specialized ATC stations.
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+ Ground control is responsible for directing all ground traffic in designated "movement areas", except the traffic on runways. This includes planes, baggage trains, snowplows, grass cutters, fuel trucks, stair trucks, airline food trucks, conveyor belt vehicles and other vehicles. Ground Control will instruct these vehicles on which taxiways to use, which runway they will use (in the case of planes), where they will park, and when it is safe to cross runways. When a plane is ready to takeoff it will be turned over to tower control. Conversely, after a plane has landed it will depart the runway and be "handed over" from Tower to Ground Control.
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+ Tower control is responsible for aircraft on the runway and in the controlled airspace immediately surrounding the airport. Tower controllers may use radar to locate an aircraft's position in 3D space, or they may rely on pilot position reports and visual observation. They coordinate the sequencing of aircraft in the traffic pattern and direct aircraft on how to safely join and leave the circuit. Aircraft which are only passing through the airspace must also contact tower control to be sure they remain clear of other traffic.
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+ At all airports the use of a traffic pattern (often called a traffic circuit outside the US) is possible. They may help to assure smooth traffic flow between departing and arriving aircraft. There is no technical need within modern commercial aviation for performing this pattern, provided there is no queue. And due to the so-called SLOT-times, the overall traffic planning tend to assure landing queues are avoided. If for instance an aircraft approaches runway 17 (which has a heading of approx. 170 degrees) from the north (coming from 360/0 degrees heading towards 180 degrees), the aircraft will land as fast as possible by just turning 10 degrees and follow the glidepath, without orbit the runway for visual reasons, whenever this is possible. For smaller piston engined airplanes at smaller airfields without ILS equipment, things are very different though.
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+
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+ Generally, this pattern is a circuit consisting of five "legs" that form a rectangle (two legs and the runway form one side, with the remaining legs forming three more sides). Each leg is named (see diagram), and ATC directs pilots on how to join and leave the circuit. Traffic patterns are flown at one specific altitude, usually 800 or 1,000 ft (244 or 305 m) above ground level (AGL). Standard traffic patterns are left-handed, meaning all turns are made to the left. One of the main reason for this is that pilots sit on the left side of the airplane, and a Left-hand patterns improves their visibility of the airport and pattern. Right-handed patterns do exist, usually because of obstacles such as a mountain, or to reduce noise for local residents. The predetermined circuit helps traffic flow smoothly because all pilots know what to expect, and helps reduce the chance of a mid-air collision.
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+ At controlled airports, a circuit can be in place but is not normally used. Rather, aircraft (usually only commercial with long routes) request approach clearance while they are still hours away from the airport; the destination airport can then plan a queue of arrivals, and planes will be guided into one queue per active runway for a "straight-in" approach. While this system keeps the airspace free and is simpler for pilots, it requires detailed knowledge of how aircraft are planning to use the airport ahead of time and is therefore only possible with large commercial airliners on pre-scheduled flights. The system has recently become so advanced that controllers can predict whether an aircraft will be delayed on landing before it even takes off; that aircraft can then be delayed on the ground, rather than wasting expensive fuel waiting in the air.
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+
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+ There are a number of aids, both visual and electronic, though not at all airports. A visual approach slope indicator (VASI) helps pilots fly the approach for landing. Some airports are equipped with a VHF omnidirectional range (VOR) to help pilots find the direction to the airport. VORs are often accompanied by a distance measuring equipment (DME) to determine the distance to the VOR. VORs are also located off airports, where they serve to provide airways for aircraft to navigate upon. In poor weather, pilots will use an instrument landing system (ILS) to find the runway and fly the correct approach, even if they cannot see the ground. The number of instrument approaches based on the use of the Global Positioning System (GPS) is rapidly increasing and may eventually become the primary means for instrument landings.
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+ Larger airports sometimes offer precision approach radar (PAR), but these systems are more common at military air bases than civilian airports. The aircraft's horizontal and vertical movement is tracked via radar, and the controller tells the pilot his position relative to the approach slope. Once the pilots can see the runway lights, they may continue with a visual landing.
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+ Airport guidance signs provide direction and information to taxiing aircraft and airport vehicles. Smaller aerodromes may have few or no signs, relying instead on diagrams and charts.
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+ Many airports have lighting that help guide planes using the runways and taxiways at night or in rain or fog.
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+
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+ On runways, green lights indicate the beginning of the runway for landing, while red lights indicate the end of the runway. Runway edge lighting consists of white lights spaced out on both sides of the runway, indicating the edges. Some airports have more complicated lighting on the runways including lights that run down the centerline of the runway and lights that help indicate the approach (an approach lighting system, or ALS). Low-traffic airports may use pilot-controlled lighting to save electricity and staffing costs.
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+ Along taxiways, blue lights indicate the taxiway's edge, and some airports have embedded green lights that indicate the centerline.
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+
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+ Weather observations at the airport are crucial to safe takeoffs and landings. In the US and Canada, the vast majority of airports, large and small, will either have some form of automated airport weather station, whether an AWOS, ASOS, or AWSS, a human observer or a combination of the two. These weather observations, predominantly in the METAR format, are available over the radio, through automatic terminal information service (ATIS), via the ATC or the flight service station.
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+
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+ Planes take-off and land into the wind to achieve maximum performance. Because pilots need instantaneous information during landing, a windsock can also be kept in view of the runway. Aviation windsocks are made with lightweight material, withstand strong winds and some are lit up after dark or in foggy weather. Because visibility of windsocks is limited, often multiple glow-orange windsocks are placed on both sides of the runway.[27]
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+ Most airports have groundcrew handling the loading and unloading of passengers, crew, baggage and other services.[citation needed] Some groundcrew are linked to specific airlines operating at the airport.
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+ Among the vehicles that serve an airliner on the ground are:
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+
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+ The length of time an aircraft remains on the ground in between consecutive flights is known as "turnaround time". Airlines pay great attention to minimizing turnaround times in an effort to keep aircraft use (flying time) high, with times scheduled as low as 25 minutes for jet aircraft operated by low-cost carriers on narrow-body aircraft.
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+ Like industrial equipment or facility management, airports require tailor-made maintenance management due to their complexity. With many tangible assets spread over a large area in different environments, these infrastructures must therefore effectively monitor these assets and store spare parts to maintain them at an optimal level of service.[28]
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+ To manage these airport assets, several solutions are competing for the market: CMMS (computerized maintenance management system) predominate, and mainly enable a company's maintenance activity to be monitored, planned, recorded and rationalized.[28]
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+ Aviation safety is an important concern in the operation of an airport, and almost every airfield includes equipment and procedures for handling emergency situations. Airport crash tender crews are equipped for dealing with airfield accidents, crew and passenger extractions, and the hazards of highly flammable aviation fuel. The crews are also trained to deal with situations such as bomb threats, hijacking, and terrorist activities.
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+ Hazards to aircraft include debris, nesting birds, and reduced friction levels due to environmental conditions such as ice, snow, or rain. Part of runway maintenance is airfield rubber removal which helps maintain friction levels. The fields must be kept clear of debris using cleaning equipment so that loose material does not become a projectile and enter an engine duct (see foreign object damage). In adverse weather conditions, ice and snow clearing equipment can be used to improve traction on the landing strip. For waiting aircraft, equipment is used to spray special deicing fluids on the wings.
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+ Many airports are built near open fields or wetlands. These tend to attract bird populations, which can pose a hazard to aircraft in the form of bird strikes. Airport crews often need to discourage birds from taking up residence.
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+ Some airports are located next to parks, golf courses, or other low-density uses of land. Other airports are located near densely populated urban or suburban areas.
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+ An airport can have areas where collisions between aircraft on the ground tend to occur. Records are kept of any incursions where aircraft or vehicles are in an inappropriate location, allowing these "hot spots" to be identified. These locations then undergo special attention by transportation authorities (such as the FAA in the US) and airport administrators.
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+
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+ During the 1980s, a phenomenon known as microburst became a growing concern due to aircraft accidents caused by microburst wind shear, such as Delta Air Lines Flight 191. Microburst radar was developed as an aid to safety during landing, giving two to five minutes' warning to aircraft in the vicinity of the field of a microburst event.
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+ Some airfields now have a special surface known as soft concrete at the end of the runway (stopway or blastpad) that behaves somewhat like styrofoam, bringing the plane to a relatively rapid halt as the material disintegrates. These surfaces are useful when the runway is located next to a body of water or other hazard, and prevent the planes from overrunning the end of the field.
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+ Airports often have on-site firefighters to respond to emergencies. These use specialized vehicles, known as airport crash tenders.
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+ Aircraft noise is a major cause of noise disturbance to residents living near airports. Sleep can be affected if the airports operate night and early morning flights. Aircraft noise occurs not only from take-offs and landings, but also from ground operations including maintenance and testing of aircraft. Noise can have other health effects as well. Other noise and environmental concerns are vehicle traffic causing noise and pollution on roads leading to airport. [29]
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+
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+ The construction of new airports or addition of runways to existing airports, is often resisted by local residents because of the effect on countryside, historical sites, and local flora and fauna. Due to the risk of collision between birds and aircraft, large airports undertake population control programs where they frighten or shoot birds.[citation needed]
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+ The construction of airports has been known to change local weather patterns. For example, because they often flatten out large areas, they can be susceptible to fog in areas where fog rarely forms. In addition, they generally replace trees and grass with pavement, they often change drainage patterns in agricultural areas, leading to more flooding, run-off and erosion in the surrounding land.[30][citation needed]
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+
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+ Some of the airport administrations prepare and publish annual environmental reports to show how they consider these environmental concerns in airport management issues and how they protect environment from airport operations. These reports contain all environmental protection measures performed by airport administration in terms of water, air, soil and noise pollution, resource conservation and protection of natural life around the airport.
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+
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+ A 2019 report from the Cooperative Research Programs of the US Transportation Research Board showed all airports have a role to play in advancing greenhouse gas (GHG) reduction initiatives. Small airports have demonstrated leadership by using their less complex organizational structure to implement newer technologies and to serve as a providing ground for their feasibility. Large airports have the economic stability and staff resources necessary to grow in-house expertise and fund comprehensive new programs.[31]
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+ A growing number of airports are installing solar photovoltaic arrays to offset their electricity use.[32][33] The National Renewable Energy Lab has shown this can be done safely.[34]
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+
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+ The world's first airport to be fully powered by solar energy is located at Kochi, India. Another airport known for considering environmental concerns is Seymour Airport in the Galapagos Islands.
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+
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+ An airbase, sometimes referred to as an air station or airfield, provides basing and support of military aircraft. Some airbases, known as military airports, provide facilities similar to their civilian counterparts. For example, RAF Brize Norton in the UK has a terminal which caters to passengers for the Royal Air Force's scheduled flights to the Falkland Islands. Some airbases are co-located with civilian airports, sharing the same ATC facilities, runways, taxiways and emergency services, but with separate terminals, parking areas and hangars. Bardufoss Airport, Bardufoss Air Station in Norway and Pune Airport in India are examples of this.
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+ An aircraft carrier is a warship that functions as a mobile airbase. Aircraft carriers allow a naval force to project air power without having to depend on local bases for land-based aircraft. After their development in World War I, aircraft carriers replaced the battleship as the centrepiece of a modern fleet during World War II.
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+ Most airports in the United States are designated "private-use airports" meaning that, whether publicly- or privately-owned, the airport is not open or available for use by the public (although use of the airport may be made available by invitation of the owner or manager).
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+
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+ Airports are uniquely represented by their IATA airport code and ICAO airport code.
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+ Most airport names include the location. Many airport names honour a public figure, commonly a politician (e.g., Charles de Gaulle Airport, George Bush Intercontinental Airport, O.R. Tambo International Airport), a monarch (e.g. Chhatrapati Shivaji International Airport, King Shaka International Airport), a cultural leader (e.g. Liverpool John Lennon Airport, Leonardo da Vinci-Fiumicino Airport, Louis Armstrong New Orleans International Airport) or a prominent figure in aviation history of the region (e.g. Sydney Kingsford Smith Airport), sometimes even famous writers (e.g. Allama Iqbal International Airport) and explorers (e.g. Venice Marco Polo Airport).
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+ Some airports have unofficial names, possibly so widely circulated that its official name is little used or even known.[citation needed]
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+ Some airport names include the word "International" to indicate their ability to handle international air traffic. This includes some airports that do not have scheduled international airline services (e.g. Port Elizabeth International Airport).
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+ The earliest aircraft takeoff and landing sites were grassy fields.[35] The plane could approach at any angle that provided a favorable wind direction. A slight improvement was the dirt-only field, which eliminated the drag from grass. However, these functioned well only in dry conditions. Later, concrete surfaces would allow landings regardless of meteorological conditions.
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+ The title of "world's oldest airport" is disputed. College Park Airport in Maryland, US, established in 1909 by Wilbur Wright, is generally agreed to be the world's oldest continuously operating airfield,[36] although it serves only general aviation traffic.
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+ Beijing Nanyuan Airport in China, which was built to accommodate planes in 1904, and airships in 1907, opened in 1910.[37] It was in operation until September 2019. Pearson Field Airport in Vancouver, Washington, United States, was built to accommodate planes in 1905 and airships in 1911, and is still in use as of January 2020.[citation needed]
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+
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+ Hamburg Airport opened in January 1911, making it the oldest commercial airport in the world which is still in operation. Bremen Airport opened in 1913 and remains in use, although it served as an American military field between 1945 and 1949. Amsterdam Airport Schiphol opened on September 16, 1916, as a military airfield, but has accepted civil aircraft only since December 17, 1920, allowing Sydney Airport—which started operations in January 1920—to claim to be one of the world's oldest continuously operating commercial airports.[38] Minneapolis-Saint Paul International Airport in the US opened in 1920 and has been in continuous commercial service since. It serves about 35,000,000 passengers each year and continues to expand, recently opening a new 11,000-foot (3,355 m) runway. Of the airports constructed during this early period in aviation, it is one of the largest and busiest that is still currently operating. Rome Ciampino Airport, opened 1916, is also a contender, as well as the Don Mueang International Airport near Bangkok, Thailand, which opened in 1914.
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+ Increased aircraft traffic during World War I led to the construction of landing fields. Aircraft had to approach these from certain directions and this led to the development of aids for directing the approach and landing slope.
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+ Following the war, some of these military airfields added civil facilities for handling passenger traffic. One of the earliest such fields was Paris – Le Bourget Airport at Le Bourget, near Paris. The first airport to operate scheduled international commercial services was Hounslow Heath Aerodrome in August 1919, but it was closed and supplanted by Croydon Airport in March 1920.[39] In 1922, the first permanent airport and commercial terminal solely for commercial aviation was opened at Flughafen Devau near what was then Königsberg, East Prussia. The airports of this era used a paved "apron", which permitted night flying as well as landing heavier aircraft.
188
+
189
+ The first lighting used on an airport was during the latter part of the 1920s; in the 1930s approach lighting came into use. These indicated the proper direction and angle of descent. The colours and flash intervals of these lights became standardized under the International Civil Aviation Organization (ICAO). In the 1940s, the slope-line approach system was introduced. This consisted of two rows of lights that formed a funnel indicating an aircraft's position on the glideslope. Additional lights indicated incorrect altitude and direction.
190
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191
+ After World War II, airport design became more sophisticated. Passenger buildings were being grouped together in an island, with runways arranged in groups about the terminal. This arrangement permitted expansion of the facilities. But it also meant that passengers had to travel further to reach their plane.
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+ An improvement in the landing field was the introduction of grooves in the concrete surface. These run perpendicular to the direction of the landing aircraft and serve to draw off excess rainwater that could build up in front of the plane's wheels.
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+ Airport construction boomed during the 1960s with the increase in jet aircraft traffic. Runways were extended out to 3,000 m (9,800 ft). The fields were constructed out of reinforced concrete using a slip-form machine that produces a continuous slab with no disruptions along the length. The early 1960s also saw the introduction of jet bridge systems to modern airport terminals, an innovation which eliminated outdoor passenger boarding. These systems became commonplace in the United States by the 1970s.[citation needed]
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+ The malicious use of UAVs has led to the deployment of counter unmanned air system (C-UAS) technologies such as the Aaronia AARTOS which have been installed on major international airports[40][41].
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+ Airports have played major roles in films and television programs due to their very nature as a transport and international hub, and sometimes because of distinctive architectural features of particular airports. One such example of this is The Terminal, a film about a man who becomes permanently grounded in an airport terminal and must survive only on the food and shelter provided by the airport. They are also one of the major elements in movies such as The V.I.P.s, Speed, Airplane!, Airport (1970), Die Hard 2, Soul Plane, Jackie Brown, Get Shorty, Home Alone, Liar Liar, Passenger 57, Final Destination (2000), Unaccompanied Minors, Catch Me If You Can, Rendition and The Langoliers. They have also played important parts in television series like Lost, The Amazing Race, America's Next Top Model, Cycle 10 which have significant parts of their story set within airports. In other programmes and films, airports are merely indicative of journeys, e.g. Good Will Hunting.
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+ Several computer simulation games put the player in charge of an airport. These include the Airport Tycoon series, SimAirport and Airport CEO.
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+ Each national aviation authority has a source of information about airports in their country. This will contain information on airport elevation, airport lighting, runway information, communications facilities and frequencies, hours of operation, nearby NAVAIDs and contact information where prior arrangement for landing is necessary.
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+ Infraero is responsible for the airports in Brazil
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+ Lists:
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+ Basque Country may refer to:
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+ A troll is a being in Scandinavian folklore, including Norse mythology. In Old Norse sources, beings described as trolls dwell in isolated rocks, mountains, or caves, live together in small family units, and are rarely helpful to human beings.
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+
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+ In later Scandinavian folklore, trolls became beings in their own right, where they live far from human habitation, are not Christianized, and are considered dangerous to human beings. Depending on the source, their appearance varies greatly; trolls may be ugly and slow-witted, or look and behave exactly like human beings, with no particularly grotesque characteristic about them.
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+
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+ Trolls are sometimes associated with particular landmarks in Scandinavian folklore, which at times may be explained as formed from a troll exposed to sunlight. Trolls are depicted in a variety of media in modern popular culture.
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+
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+ The Old Norse nouns troll and tröll (variously meaning "fiend, demon, werewolf, jötunn") and Middle High German troll, trolle "fiend" (according to philologist Vladimir Orel likely borrowed from Old Norse) developed from Proto-Germanic neuter noun *trullan. The origin of the Proto-Germanic word is unknown.[1] Additionally, the Old Norse verb trylla 'to enchant, to turn into a troll' and the Middle High German verb trüllen "to flutter" both developed from the Proto-Germanic verb *trulljanan, a derivative of *trullan.[1]
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+
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+ In Norse mythology, troll, like thurs, is a term applied to jötnar and is mentioned throughout the Old Norse corpus. In Old Norse sources, trolls are said to dwell in isolated mountains, rocks, and caves, sometimes live together (usually as father-and-daughter or mother-and-son), and are rarely described as helpful or friendly.[2] The Prose Edda book Skáldskaparmál describes an encounter between an unnamed troll woman and the 9th-century skald Bragi Boddason. According to the section, Bragi was driving through "a certain forest" late one evening when a troll woman aggressively asked him who he was, in the process describing herself:
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+
13
+ Old Norse:
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+
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+ Anthony Faulkes translation:
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+
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+ John Lindow translation:
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+
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+ Bragi responds in turn, describing himself and his abilities as a skillful skald, before the scenario ends.[4]
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+
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+ There is much confusion and overlap in the use of Old Norse terms jötunn, troll, þurs, and risi, which describe various beings. Lotte Motz theorized that these were originally four distinct classes of beings: lords of nature (jötunn), mythical magicians (troll), hostile monsters (þurs), and heroic and courtly beings (risi), the last class being the youngest addition. On the other hand, Ármann Jakobson is critical of Motz's interpretation and calls this theory "unsupported by any convincing evidence".[5] Ármann highlights that the term is used to denote various beings, such as a jötunn or mountain-dweller, a witch, an abnormally strong or large or ugly person, an evil spirit, a ghost, a blámaðr, a magical boar, a heathen demi-god, a demon, a brunnmigi, or a berserker.[6]
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+ Later in Scandinavian folklore, trolls become defined as a particular type of being.[7] Numerous tales are recorded about trolls in which they are frequently described as being extremely old, very strong, but slow and dim-witted, and are at times described as man-eaters and as turning to stone upon contact with sunlight.[8] However, trolls are also attested as looking much the same as human beings, without any particularly hideous appearance about them, but living far away from human habitation and generally having "some form of social organization"—unlike the rå and näck, who are attested as "solitary beings". According to John Lindow, what sets them apart is that they are not Christian, and those who encounter them do not know them. Therefore, trolls were in the end dangerous, regardless of how well they might get along with Christian society, and trolls display a habit of bergtagning ('kidnapping'; literally "mountain-taking") and overrunning a farm or estate.[9]
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+ Lindow states that the etymology of the word "troll" remains uncertain, though he defines trolls in later Swedish folklore as "nature beings" and as "all-purpose otherworldly being[s], equivalent, for example, to fairies in Anglo-Celtic traditions". They "therefore appear in various migratory legends where collective nature-beings are called for". Lindow notes that trolls are sometimes swapped out for cats and "little people" in the folklore record.[9]
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+ A Scandinavian folk belief that lightning frightens away trolls and jötnar appears in numerous Scandinavian folktales, and may be a late reflection of the god Thor's role in fighting such beings. In connection, the lack of trolls and jötnar in modern Scandinavia is sometimes explained as a result of the "accuracy and efficiency of the lightning strokes".[10] Additionally, the absence of trolls in regions of Scandinavia is described in folklore as being a "consequence of the constant din of the church-bells". This ring caused the trolls to leave for other lands, although not without some resistance; numerous traditions relate how trolls destroyed a church under construction or hurled boulders and stones at completed churches. Large local stones are sometimes described as the product of a troll's toss.[11] Additionally, into the 20th century, the origins of particular Scandinavian landmarks, such as particular stones, are ascribed to trolls who may, for example, have turned to stone upon exposure to sunlight.[8]
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+
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+ Lindow compares the trolls of the Swedish folk tradition to Grendel, the supernatural mead hall invader in the Old English poem Beowulf, and notes that "just as the poem Beowulf emphasizes not the harrying of Grendel but the cleansing of the hall of Beowulf, so the modern tales stress the moment when the trolls are driven off."[9]
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+ Smaller trolls are attested as living in burial mounds and in mountains in Scandinavian folk tradition.[12] In Denmark, these creatures are recorded as troldfolk ("troll-folk"), bjergtrolde ("mountain-trolls"), or bjergfolk ("mountain-folk") and in Norway also as troldfolk ("troll-folk") and tusser.[12] Trolls may be described as small, human-like beings or as tall as men depending on the region of origin of the story.[13]
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+ In Norwegian tradition, similar tales may be told about the larger trolls and the Huldrefolk ("hidden-folk"), yet a distinction is made between the two. The use of the word trow in Orkney and Shetland, to mean beings which are very like the Huldrefolk in Norway, may suggest a common origin for the terms. The word troll may have been used by pagan Norse settlers in Orkney and Shetland as a collective term for supernatural beings who should be respected and avoided rather than worshipped. Troll could later have become specialized as a description of the larger, more menacing Jötunn-kind whereas Huldrefolk may have developed as the term for smaller trolls.[14]
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+ John Arnott MacCulloch posited a connection between the Old Norse vættir and trolls, suggesting that both concepts may derive from spirits of the dead.[15]
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+ Troll, a Norwegian research station in Antarctica, is so named because of the rugged mountains which stand around that place like trolls. It includes a ground station which tracks satellites in polar orbit.
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+ Media related to trolls at Wikimedia Commons
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+ The Fallopian tubes, also known as uterine tubes or salpinges (singular salpinx), are tubes that stretch from the uterus to the ovaries, and are part of the female reproductive system.
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+ The fertilized egg passes through the Fallopian tubes from the ovaries of female mammals to the uterus. The Fallopian tubes is simple columnar epithelium with hair-like extensions called cilia which carry the fertilized egg. In non-mammalian vertebrates, the equivalent of a Fallopian tube is an oviduct.
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+ The name comes from the Catholic priest and anatomist Gabriele Falloppio, for whom other anatomical structures are also named.
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+
7
+ The Fallopian tube is composed of four parts. These are, described from near the ovaries to inwards near the uterus, the infundibulum with its associated fimbriae near the ovary, the ampulla that represents the major portion of the lateral tube, the isthmus, which is the narrower part of the tube that links to the uterus, and the interstitial (or intramural) part, the narrowest part of the uterine tube, that crosses the muscles of the uterine. The average length of a fallopian tube is 11-12 cm.[1]
8
+
9
+ The uterus opens into the Fallopian tube at the proximal tubal opening (also called the proximal ostium or os[2]), after the uterotubal junction, and accessible via hysteroscopy. Occlusion at this opening is referred to as proximal tubal occlusion. From there there are three named parts of the Fallopian tube; the isthmus, the ampulla, and the infundibulum. The isthmus sits next to the opening of the Fallopian tube into the uterus. It connects to the ampulla (Latin: flask), which curves over the ovary and is the most common site of human fertilization.
10
+
11
+ The ampulla connects with the infundibulum, which rests above the ovaries, and ends at the distal tubal opening (or abdominal ostium)[3] into the abdominal cavity where, in ovulation, the oocyte enters the Fallopian tube. The opening is surrounded by fimbriae, which help in the collection of the oocyte. Occlusion of this opening is referred to as distal tubal occlusion. The fimbriae (singular fimbria) is a fringe of tissue around the ostium of the Fallopian tube, in the direction of the ovary. Of all fimbriae, one fimbria is long enough to reach the ovary. It is called fimbria ovarica.[4][5]
12
+
13
+ An ovary is not directly connected to its adjacent Fallopian tube. When ovulation is about to occur, the sex hormones activate the fimbriae, causing them to swell with blood and hit the ovary in a gentle, sweeping motion. An oocyte is released from the ovary into the peritoneal cavity and the cilia of the fimbriae sweep the ovum into the Fallopian tube.
14
+
15
+ When viewed under the microscope, the Fallopian tube has four layers. From outer to inner these are the serosa, subserosa, lamina propria and innermost mucosa. The serosa is derived from visceral peritoneum. The subserosa is composed of loose adventitious tissue, blood vessels, lymphatics, an outer longitudinal and inner circular smooth muscle coats. This layer is responsible for the rhythmic contraction, called peristalsis, of the Fallopian tubes. Lamina propria is a vascular connective tissue.[6] The histological features of tube vary along its length. The mucosa of the ampulla contains an extensive array of complex folds, whereas the relatively narrow isthmus has a thick muscular coat and simple mucosal folds.[6]
16
+
17
+ The innermost layer of the tube is a single layer of column-shaped cells. The columnar cells have microscopic hair-like filaments called cilia throughout the tube, most numerous in the infundibulum and ampulla. Estrogen increases the production of cilia on these cells. Between the ciliated cells are peg cells, which contain apical granules and produce tubular fluid. This fluid contains nutrients for spermatozoa, oocytes, and zygotes. The secretions also promote capacitation of the sperm by removing glycoproteins and other molecules from the plasma membrane of the sperm. Progesterone increases the number of peg cells, while estrogen increases their height and secretory activity. Fluid flows through the tubes towards the ovaries, the opposite direction to the action of the cilia.
18
+
19
+ Embryos develop a urogenital ridge that forms at their tail end and eventually forms the basis for the urinary system and reproductive tracts. Either side and to the front of this tract, around the sixth week develops a duct called the paramesonephric duct, also called the Müllerian duct.[7] A second duct, the mesonephric duct, develops adjacent to this. Both ducts become longer over the next two weeks, and the paramesonephric ducts around the eighth week cross to meet in the midline and fuse.[7] One duct then regresses, with this depending on whether the embryo is genetically female or male. In females, the paramesonephric duct remains, and eventually forms the female reproductive tract.[7] The portions of the paramesonephric duct which are more cranial - that is, further from the tail-end, end up forming the fallopian tubes.[7] In males, because of the presence of the Y sex chromosome, anti-mullerian hormone is produced. This leads to the degeneration of the paramesonephric duct.[7]
20
+
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+ As the uterus develops, the part of the fallopian tubes closer to the uterus, the ampulla, become larger. Extensions from the fallopian tubes, the fimbriae, develop over time.
22
+
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+ Apart from the presence of sex chromosomes, specific genes associated with the development of the fallopian tubes include the Wnt and Hox groups of genes, Lim1, Pax2, and Emx2.[7]
24
+
25
+ Embryos have two pairs of ducts to let gametes out of the body; one pair (the Müllerian ducts) develops in females into the fallopian tubes, uterus and vagina, while the other pair (the Wolffian ducts) develops in males into the epididymis and vas deferens.
26
+
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+ The homologous organ in the male is the rudimentary appendix testis.[citation needed]
28
+
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+ The fallopian tube allows passage of an egg from the ovary to the uterus. When an oocyte is developing in an ovary, it is surrounded by a spherical collection of cells known as an ovarian follicle. Just before ovulation, the primary oocyte completes meiosis I to form the first polar body and a secondary oocyte which is arrested in metaphase of meiosis II.
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+
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+ At the time of ovulation in the menstrual cycle, the secondary oocyte is released from the ovary. The follicle and the ovary's wall rupture, allowing the secondary oocyte to escape. The secondary oocyte is caught by the fimbriated end of the Fallopian tube and travels to the ampulla. Here, the egg is able to become fertilised with sperm. The ampulla is typically where the sperm are met and fertilization occurs; meiosis II is promptly completed. After fertilisation, the ovum is now called a zygote and travels towards the uterus with the aid of the hair-like cilia and the activity of the muscle of the Fallopian tube. The early embryo requires critical development in the fallopian tube.[8] After about five days the new embryo enters the uterine cavity and on about the sixth day implants on the wall of the uterus.
32
+
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+ The release of an oocyte does not alternate between the two ovaries and seems to be random. After removal of an ovary, the remaining one produces an egg every month.[9]
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+
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+ Salpingitis is inflammation of the fallopian tubes and may be found alone, or as part of pelvic inflammatory disease (PID). A thickening of the fallopian tube at its narrow portion, due to inflammation, is known as salpingitis isthmica nodosa. Like PID and endometriosis, it may lead to fallopian tube obstruction. Fallopian tube obstruction may be a cause of infertility or ectopic pregnancy.[10]
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+
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+ Fallopian tube cancer, which typically arises from the epithelial lining of the fallopian tube, has historically been considered to be a very rare malignancy. Recent evidence suggests it probably represents a significant portion of what has been classified as ovarian cancer in the past.[11] While tubal cancers may be misdiagnosed as ovarian cancer, it is of little consequence as the treatment of both ovarian and fallopian tube cancer is similar.[citation needed]
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+
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+ Occasionally the embryo implants into the fallopian tube instead of the uterus, creating an ectopic pregnancy, commonly known as a "tubal pregnancy".
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+
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+ While a full testing of tubal functions in patients with infertility is not possible, testing of whether the tubes are open, called patency, is important as tubal obstruction is a major cause of infertility. A hysterosalpingogram, laparoscopy and dye, or hysterocontrast sonography will demonstrate whether the tubes are open. Tubal insufflation is a standard procedure for testing patency. During surgery the condition of the tubes may be inspected and a dye such as methylene blue can be injected into the uterus and shown to pass through the tubes when the cervix is occluded. As tubal disease is often related to Chlamydia infection, testing for Chlamydia antibodies has become a cost-effective screening device for tubal pathology.[12]
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+
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+ The surgical removal of a fallopian tube is called a salpingectomy. To remove both tubes is a bilateral salpingectomy. An operation that combines the removal of a fallopian tube with removal of at least one ovary is a salpingo-oophorectomy. An operation to remove a fallopian tube obstruction is called a tuboplasty.
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+
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+ The fallopian tube can prolapse into the vagina and can be mistaken for a tumour. When this happens, it is usually after a hysterectomy.[13]
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+
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+ The Fallopian tubes are named after the 16th-century Italian anatomist Gabriele Falloppio, the first person to provide a detailed description of the tubes.[14][15] He thought they resembled tubas, the plural of tuba in Italian being tube which was misunderstood and became the English "tube".[16]
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+ Though the name Fallopian tube is eponymous, it is often spelt with a lower case f from the assumption that the adjective fallopian has been absorbed into modern English as the de facto name for the structure.[citation needed]
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+ Image showing the right Fallopian tube (here labelled the uterine tube) seen from behind. The uterus, ovaries and right broad ligament are labelled.
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+
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+ Unlabelled image showing the right Fallopian tube
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+
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+ Isthmus of the Fallopian tube seen arising from the uterus in a cadaveric specimen.
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+ This article incorporates text in the public domain from page 1257 of the 20th edition of Gray's Anatomy (1918)
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1
+ Brass
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+
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+ Plucked
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+ The trumpet is a brass instrument commonly used in classical and jazz ensembles. The trumpet group ranges from the piccolo trumpet with the highest register in the brass family, to the bass trumpet, which is pitched one octave below the standard B♭ or C Trumpet.
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+
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+ Trumpet-like instruments have historically been used as signalling devices in battle or hunting, with examples dating back to at least 1500 BC. They began to be used as musical instruments only in the late 14th or early 15th century.[1] Trumpets are used in art music styles, for instance in orchestras, concert bands, and jazz ensembles, as well as in popular music. They are played by blowing air through nearly-closed lips (called the player's embouchure), producing a "buzzing" sound that starts a standing wave vibration in the air column inside the instrument.[2] Since the late 15th century, trumpets have primarily been constructed of brass tubing, usually bent twice into a rounded rectangular shape.[3]
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+
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+ There are many distinct types of trumpet, with the most common being pitched in B♭ (a transposing instrument), having a tubing length of about 1.48 m (4 ft 10 in). Early trumpets did not provide means to change the length of tubing, whereas modern instruments generally have three (or sometimes four) valves in order to change their pitch. There are eight combinations of three valves, making seven different tubing lengths, with the third valve sometimes used as an alternate fingering equivalent to the 1–2 combination. Most trumpets have valves of the piston type, while some have the rotary type. The use of rotary-valved trumpets is more common in orchestral settings (especially in German and German-style orchestras), although this practice varies by country. Each valve, when engaged, increases the length of tubing, lowering the pitch of the instrument. A musician who plays the trumpet is called a trumpet player or trumpeter.[4]
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+
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+ The English word "trumpet" was first used in the late 14th century.[5] The word came from Old French "trompette", which is a diminutive of trompe.[5] The word "trump", meaning "trumpet," was first used in English in 1300. The word comes from Old French trompe "long, tube-like musical wind instrument" (12c.), cognate with Provençal tromba, Italian tromba, all probably from a Germanic source (compare Old High German trumpa, Old Norse trumba "trumpet"), of imitative origin."[6]
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+
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+ The earliest trumpets date back to 1500 BC and earlier. The bronze and silver trumpets from Tutankhamun's grave in Egypt, bronze lurs from Scandinavia, and metal trumpets from China date back to this period.[7] Trumpets from the Oxus civilization (3rd millennium BC) of Central Asia have decorated swellings in the middle, yet are made out of one sheet of metal, which is considered a technical wonder.[8]
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+
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+ The Shofar, made from a ram horn and the Hatzotzeroth, made of metal, are both mentioned in the Bible. They were played in Solomon's Temple around 3000 years ago. They were said to be used to blow down the walls of Jericho. They are still used on certain religious days.[9]
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+ The Salpinx was a straight trumpet 62 inches (1,600 mm) long, made of bone or bronze. Salpinx contests were a part of the original Olympic Games.[9]
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+
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+ The Moche people of ancient Peru depicted trumpets in their art going back to AD 300.[10] The earliest trumpets were signaling instruments used for military or religious purposes, rather than music in the modern sense;[11] and the modern bugle continues this signaling tradition.
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+
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+ Improvements to instrument design and metal making in the late Middle Ages and Renaissance led to an increased usefulness of the trumpet as a musical instrument. The natural trumpets of this era consisted of a single coiled tube without valves and therefore could only produce the notes of a single overtone series. Changing keys required the player to change crooks of the instrument.[9] The development of the upper, "clarino" register by specialist trumpeters—notably Cesare Bendinelli—would lend itself well to the Baroque era, also known as the "Golden Age of the natural trumpet." During this period, a vast body of music was written for virtuoso trumpeters. The art was revived in the mid-20th century and natural trumpet playing is again a thriving art around the world. Many modern players in Germany and the UK who perform Baroque music use a version of the natural trumpet fitted with three or four vent holes to aid in correcting out-of-tune notes in the harmonic series.[12]
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+
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+ The melody-dominated homophony of the classical and romantic periods relegated the trumpet to a secondary role by most major composers owing to the limitations of the natural trumpet. Berlioz wrote in 1844:
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+
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+ Notwithstanding the real loftiness and distinguished nature of its quality of tone, there are few instruments that have been more degraded (than the trumpet). Down to Beethoven and Weber, every composer – not excepting Mozart – persisted in confining it to the unworthy function of filling up, or in causing it to sound two or three commonplace rhythmical formulae.[13]
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+
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+ The trumpet is constructed of brass tubing bent twice into a rounded oblong shape.[14] As with all brass instruments, sound is produced by blowing air through closed lips, producing a "buzzing" sound into the mouthpiece and starting a standing wave vibration in the air column inside the trumpet. The player can select the pitch from a range of overtones or harmonics by changing the lip aperture and tension (known as the embouchure).
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+
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+ The mouthpiece has a circular rim, which provides a comfortable environment for the lips' vibration. Directly behind the rim is the cup, which channels the air into a much smaller opening (the back bore or shank) that tapers out slightly to match the diameter of the trumpet's lead pipe. The dimensions of these parts of the mouthpiece affect the timbre or quality of sound, the ease of playability, and player comfort. Generally, the wider and deeper the cup, the darker the sound and timbre.
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+
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+ Modern trumpets have three (or, infrequently, four) piston valves, each of which increases the length of tubing when engaged, thereby lowering the pitch. The first valve lowers the instrument's pitch by a whole step (two semitones), the second valve by a half step (one semitone), and the third valve by one and a half steps (three semitones). When a fourth valve is present, as with some piccolo trumpets, it usually lowers the pitch a perfect fourth (five semitones). Used singly and in combination these valves make the instrument fully chromatic, i.e., able to play all twelve pitches of classical music. For more information about the different types of valves, see Brass instrument valves.
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+
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+ The pitch of the trumpet can be raised or lowered by the use of the tuning slide. Pulling the slide out lowers the pitch; pushing the slide in raises it. To overcome the problems of intonation and reduce the use of the slide, Renold Schilke designed the tuning-bell trumpet. Removing the usual brace between the bell and a valve body allows the use of a sliding bell; the player may then tune the horn with the bell while leaving the slide pushed in, or nearly so, thereby improving intonation and overall response.[15]
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+
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+ A trumpet becomes a closed tube when the player presses it to the lips; therefore, the instrument only naturally produces every other overtone of the harmonic series. The shape of the bell makes the missing overtones audible.[16] Most notes in the series are slightly out of tune and modern trumpets have slide mechanisms for the first and third valves with which the player can compensate by throwing (extending) or retracting one or both slides, using the left thumb and ring finger for the first and third valve slides respectively.
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+
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+ The most common type is the B♭ trumpet, but A, C, D, E♭, E, low F, and G trumpets are also available. The C trumpet is most common in American orchestral playing, where it is used alongside the B♭ trumpet. Orchestral trumpet players are adept at transposing music at sight, frequently playing music written for the A, B♭, D, E♭, E, or F trumpet on the C trumpet or B♭ trumpet.
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+
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+ The smallest trumpets are referred to as piccolo trumpets. The most common of these are built to play in both B♭ and A, with separate leadpipes for each key. The tubing in the B♭ piccolo trumpet is one-half the length of that in a standard B♭ trumpet. Piccolo trumpets in G, F and C are also manufactured, but are less common. Many players use a smaller mouthpiece on the piccolo trumpet, which requires a different sound production technique from the B♭ trumpet and can limit endurance. Almost all piccolo trumpets have four valves instead of the usual three — the fourth valve lowers the pitch, usually by a fourth, to assist in the playing of lower notes and to create alternate fingerings that facilitate certain trills. Maurice André, Håkan Hardenberger, David Mason, and Wynton Marsalis are some well-known trumpet players known for their additional virtuosity on the piccolo trumpet.
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+
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+ Trumpets pitched in the key of low G are also called sopranos, or soprano bugles, after their adaptation from military bugles. Traditionally used in drum and bugle corps, sopranos have featured both rotary valves and piston valves.
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+
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+ The bass trumpet is usually played by a trombone player, being at the same pitch.[4] Bass trumpet is played with a shallower trombone mouthpiece, and music for it is written in treble clef. The most common keys for bass trumpets are C and B♭. Both C and B♭ bass trumpets are transposing instruments sounding an octave (C) or a major ninth (B♭) lower than written.
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+
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+ The historical slide trumpet was probably first developed in the late 14th century for use in alta cappella wind bands. Deriving from early straight trumpets, the Renaissance slide trumpet was essentially a natural trumpet with a sliding leadpipe. This single slide was rather awkward, as the entire corpus of the instrument moved, and the range of the slide was probably no more than a major third. Originals were probably pitched in D, to fit with shawms in D and G, probably at a typical pitch standard near A=466 Hz. As no known instruments from this period survive, the details—and even the existence—of a Renaissance slide trumpet is a matter of conjecture and debate among scholars.[17]
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+
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+ Some slide trumpet designs saw use in England in the 18th century.[18]
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+
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+ The pocket trumpet is a compact B♭ trumpet. The bell is usually smaller than a standard trumpet and the tubing is more tightly wound to reduce the instrument size without reducing the total tube length. Its design is not standardized, and the quality of various models varies greatly. It can have a tone quality and projection unique in the trumpet world: a warm sound and a voice-like articulation. Since many pocket trumpet models suffer from poor design as well as cheap and imprecise manufacturing, the intonation, tone color and dynamic range of such instruments are severely hindered. Professional-standard instruments are, however, available. While they are not a substitute for the full-sized instrument, they can be useful in certain contexts. The jazz musician Don Cherry was renowned for his playing of the pocket instrument.
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+
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+ The herald trumpet has an elongated bell extending far in front of the player, allowing a standard length of tubing from which a flag may be hung; the instrument is mostly used for ceremonial events such as parades and fanfares.
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+
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+ Monette designed the flumpet in 1989 for jazz musician Art Farmer. It is a hybrid instrument with elements of trumpet and flugelhorn, sharing the three piston valve design and with a pitch of B♭.[19]
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+ There are also rotary-valve, or German, trumpets (which are commonly used in professional German and Austrian orchestras) as well as alto and Baroque trumpets.
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+
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+ Another variant of the standard trumpet is the Vienna valve trumpet. Primarily used in Viennese brass ensembles and orchestras such as the Vienna Philharmonic and Mnozil Brass.
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+ The trumpet is often confused with its close relative the cornet, which has a more conical tubing shape compared to the trumpet's more cylindrical tube. This, along with additional bends in the cornet's tubing, gives the cornet a slightly mellower tone, but the instruments are otherwise nearly identical. They have the same length of tubing and, therefore, the same pitch, so music written for cornet and trumpet is interchangeable. Another relative, the flugelhorn, has tubing that is even more conical than that of the cornet, and an even richer tone. It is sometimes augmented with a fourth valve to improve the intonation of some lower notes.
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+
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+ On any modern trumpet, cornet, or flugelhorn, pressing the valves indicated by the numbers below produces the written notes shown. "Open" means all valves up, "1" means first valve, "1–2" means first and second valve simultaneously, and so on. The sounding pitch depends on the transposition of the instrument. Engaging the fourth valve, if present, usually drops any of these pitches by a perfect fourth as well. Within each overtone series, the different pitches are attained by changing the embouchure. Standard fingerings above high C are the same as for the notes an octave below (C♯ is 1–2, D is 1, etc.)[20]
61
+
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+ Each overtone series on the trumpet begins with the first overtone—the fundamental of each overtone series cannot be produced except as a pedal tone. Notes in parentheses are the sixth overtone, representing a pitch with a frequency of seven times that of the fundamental; while this pitch is close to the note shown, it is slightly flat relative to equal temperament, and use of those fingerings is generally avoided.
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+
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+ The fingering schema arises from the length of each valve's tubing (a longer tube produces a lower pitch). Valve "1" increases the tubing length enough to lower the pitch by one whole step, valve "2" by one half step, and valve "3" by one and a half steps. This scheme and the nature of the overtone series create the possibility of alternate fingerings for certain notes. For example, third-space "C" can be produced with no valves engaged (standard fingering) or with valves 2–3. Also, any note produced with 1–2 as its standard fingering can also be produced with valve 3 – each drops the pitch by ​1 1⁄2 steps. Alternate fingerings may be used to improve facility in certain passages, or to aid in intonation. Extending the third valve slide when using the fingerings 1–3 or 1-2-3 further lowers the pitch slightly to improve intonation.
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+ If you take the partials of the harmonic series that a modern Bb trumpet can play for each combination of valves pressed, some notes are in tune with 12-tone equal temperament and some are not. The following tables show all notes from all partials from all valve combinations in order up to high C (C6). This often gives multiple fingerings for the same note and some notes that aren't usually useful.
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+
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+ Various types of mutes can be used to alter the sound of the instrument when placed in or over the bell. While most types of mutes do decrease the volume the instrument produces, as the name implies, the sound modification is typically the primary reason for their use. Types of mutes most commonly used to alter the sound of the instrument are: Straight Mutes, Harmon Mutes (aka "Wah-Wah" Mutes), Plunger Mutes, Bucket Mutes, and Cup Mutes. A description of their construction and sound quality are below:
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+
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+ Straight Mute: Constructed of either aluminum, which produces a bright piercing sound, or stone lined with cardboard, which produces a stuffy sound.
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+
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+ Harmon Mute: Constructed of aluminum and consists of two parts called the "stem" and the "body". The stem can be extended or removed to produce different timbres of sound. This mute is also called the "Wah-Wah" mute due to its distinctive sound created by the player placing their hand over the stem opening and waving it back and forth.
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+
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+ Plunger Mute: Most often made of a rubber bathroom plunger without the stick. This is used to manipulate sound by the player holding it over the bell with their left hand.
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+ Bucket Mute: Constructed from cardboard and cloth, this mute is clipped to the end of the bell and used to muffle the sound almost completely.
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+ Cup Mute: Also constructed of cardboard, this mute is shaped exactly like a straight mute but includes a cup at the end. In many models the cup is adjustable much like the stem on the harmon mute and produces a softer more muffled sound than a traditional straight mute.
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+ The standard trumpet range extends from the written F♯ immediately below Middle C up to about three octaves higher (F♯3 – F♯6). Traditional trumpet repertoire rarely calls for notes beyond this range, and the fingering tables of most method books peak at the high C, two octaves above middle C.[contradictory] Several trumpeters have achieved fame for their proficiency in the extreme high register, among them Maynard Ferguson, Cat Anderson, Dizzy Gillespie, Doc Severinsen, and more recently Wayne Bergeron, Thomas Gansch, James Morrison, Jon Faddis and Arturo Sandoval. It is also possible to produce pedal tones below the low F♯, which is a device occasionally employed in the contemporary repertoire for the instrument.
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+ Contemporary music for the trumpet makes wide uses of extended trumpet techniques.
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+ Flutter tonguing: The trumpeter rolls the tip of the tongue to produce a 'growling like' tone. It is achieved as if one were rolling an R in the Spanish language. This technique is widely employed by composers like Berio and Stockhausen.
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+ Growling: Simultaneously playing tone while using the back of the tongue to vibrate the uvula creating a distinct sound. Most trumpet players will use a plunger with this technique to achieve a particular sound heard in a lot of Chicago Jazz of the 1950s.
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+ Double tonguing: The player articulates using the syllables ta-ka ta-ka ta-ka
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+ Triple tonguing: The same as double tonguing, but with the syllables ta-ta-ka ta-ta-ka ta-ta-ka or ta-ka-ta ta-ka-ta.
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+
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+ Doodle tongue: The trumpeter tongues as if saying the word doodle. This is a very faint tonguing similar in sound to a valve tremolo.
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+
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+ Glissando: Trumpeters can slide between notes by depressing the valves halfway and changing the lip tension. Modern repertoire makes extensive use of this technique.
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+ Vibrato: It is often regulated in contemporary repertoire through specific notation. Composers can call for everything from fast, slow or no vibrato to actual rhythmic patterns played with vibrato.
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+ Pedal tone: Composers have written for two-and-a-half octaves below the low F♯, which is at the bottom of the standard range. Extreme low pedals are produced by slipping the lower lip out of the mouthpiece. Claude Gordon assigned pedals as part of his trumpet practice routines, that were a systematic expansion on his lessons with Herbert L. Clarke. The technique was pioneered by Bohumir Kryl.[21]
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+
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+ Microtones: Composers such as Scelsi and Stockhausen have made wide use of the trumpet's ability to play microtonally. Some instruments feature a fourth valve that provides a quarter-tone step between each note. The jazz musician Ibrahim Maalouf uses such a trumpet, invented by his father to make it possible to play Arab maqams.
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+
102
+ Valve tremolo: Many notes on the trumpet can be played in several different valve combinations. By alternating between valve combinations on the same note, a tremolo effect can be created. Berio makes extended use of this technique in his Sequenza X.
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+
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+ Noises: By hissing, clicking, or breathing through the instrument, the trumpet can be made to resonate in ways that do not sound at all like a trumpet. Noises may require amplification.
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+ Preparation: Composers have called for trumpeters to play under water, or with certain slides removed. It is increasingly common for composers to specify all sorts of preparations for trumpet. Extreme preparations involve alternate constructions, such as double bells and extra valves.
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+ Split tone: Trumpeters can produce more than one tone simultaneously by vibrating the two lips at different speeds. The interval produced is usually an octave or a fifth.
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+ Lip-trill or shake: Also known as "lip-slurs". By rapidly varying air speed, but not changing the depressed valves, the pitch varies quickly between adjacent harmonic partials. Shakes and lip-trills can vary in speed and the distance between the partials can be as large or small as the musicians' desires. Traditionally, however, lip-trills and shakes are usually the next partial up from the written note.
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+ Multi-phonics: Playing a note and "humming" a different note simultaneously. For example, sustaining a middle C and humming a major 3rd "E" at the same time.
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+ Circular breathing: A technique wind players use to produce uninterrupted tone, without pauses for breaths. The player puffs up the cheeks, storing air, then breathes in rapidly through the nose while using the cheeks to continue pushing air outwards.
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+ One trumpet method is Jean-Baptiste Arban's Complete Conservatory Method for Trumpet (Cornet).[22] Other well-known method books include Technical Studies by Herbert L. Clarke,[23] Grand Method by Louis Saint-Jacome, Daily Drills and Technical Studies by Max Schlossberg, and methods by Ernest S. Williams, Claude Gordon, Charles Colin, James Stamp, and Louis Davidson.[24] A common method book for beginners is the Walter Beeler's Method for the Cornet, and there have been several instruction books written by virtuoso Allen Vizzutti.[25] Merri Franquin wrote a Complete Method for Modern Trumpet,[26] which fell into obscurity for much of the twentieth century until public endorsements by Maurice André revived interest in this work.[27]
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+
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+ In early jazz, Louis Armstrong was well known for his virtuosity and his improvisations on the Hot Five and Hot Seven recordings, and his switch from cornet to trumpet is often cited as heralding the trumpet's dominance over the cornet in jazz.[4][28] Dizzy Gillespie was a gifted improviser with an extremely high (but musical) range, building on the style of Roy Eldridge but adding new layers of harmonic complexity. Gillespie had an enormous impact on virtually every subsequent trumpeter, both by the example of his playing and as a mentor to younger musicians. Miles Davis is widely considered one of the most influential musicians of the 20th century—his style was distinctive and widely imitated. Davis' phrasing and sense of space in his solos have been models for generations of jazz musicians.[29] Cat Anderson was a trumpet player who was known for the ability to play extremely high with an even more extreme volume, who played with Duke Ellington's Big Band. Maynard Ferguson came to prominence playing in Stan Kenton's orchestra, before forming his own band in 1957. He was noted for being able to play accurately in a remarkably high register.[30]
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+
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+ Anton Weidinger developed in the 1790s the first successful keyed trumpet, capable of playing all the chromatic notes in its range. Joseph Haydn's Trumpet Concerto was written for him in 1796 and startled contemporary audiences by its novelty,[31] a fact shown off by some stepwise melodies played low in the instrument's range.
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+
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+ The Last Judgment (workshop of Hieronymus Bosch) [nl], c.1500-1510
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+
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+ Trumpet-Player in front of a Banquet, Gerrit Dou, c.1660-1665
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+
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+ Illustration for The Trumpeter Taken Prisoner from Baby's Own Aesop, a children's edition of Aesop's fables
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+ Louis Armstrong statue in Algiers, New Orleans
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+ Miles Davis statue in Kielce, Poland
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1
+ Brass
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+
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+ Plucked
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+
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+ The trumpet is a brass instrument commonly used in classical and jazz ensembles. The trumpet group ranges from the piccolo trumpet with the highest register in the brass family, to the bass trumpet, which is pitched one octave below the standard B♭ or C Trumpet.
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+ Trumpet-like instruments have historically been used as signalling devices in battle or hunting, with examples dating back to at least 1500 BC. They began to be used as musical instruments only in the late 14th or early 15th century.[1] Trumpets are used in art music styles, for instance in orchestras, concert bands, and jazz ensembles, as well as in popular music. They are played by blowing air through nearly-closed lips (called the player's embouchure), producing a "buzzing" sound that starts a standing wave vibration in the air column inside the instrument.[2] Since the late 15th century, trumpets have primarily been constructed of brass tubing, usually bent twice into a rounded rectangular shape.[3]
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+
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+ There are many distinct types of trumpet, with the most common being pitched in B♭ (a transposing instrument), having a tubing length of about 1.48 m (4 ft 10 in). Early trumpets did not provide means to change the length of tubing, whereas modern instruments generally have three (or sometimes four) valves in order to change their pitch. There are eight combinations of three valves, making seven different tubing lengths, with the third valve sometimes used as an alternate fingering equivalent to the 1–2 combination. Most trumpets have valves of the piston type, while some have the rotary type. The use of rotary-valved trumpets is more common in orchestral settings (especially in German and German-style orchestras), although this practice varies by country. Each valve, when engaged, increases the length of tubing, lowering the pitch of the instrument. A musician who plays the trumpet is called a trumpet player or trumpeter.[4]
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+
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+ The English word "trumpet" was first used in the late 14th century.[5] The word came from Old French "trompette", which is a diminutive of trompe.[5] The word "trump", meaning "trumpet," was first used in English in 1300. The word comes from Old French trompe "long, tube-like musical wind instrument" (12c.), cognate with Provençal tromba, Italian tromba, all probably from a Germanic source (compare Old High German trumpa, Old Norse trumba "trumpet"), of imitative origin."[6]
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+
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+ The earliest trumpets date back to 1500 BC and earlier. The bronze and silver trumpets from Tutankhamun's grave in Egypt, bronze lurs from Scandinavia, and metal trumpets from China date back to this period.[7] Trumpets from the Oxus civilization (3rd millennium BC) of Central Asia have decorated swellings in the middle, yet are made out of one sheet of metal, which is considered a technical wonder.[8]
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+
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+ The Shofar, made from a ram horn and the Hatzotzeroth, made of metal, are both mentioned in the Bible. They were played in Solomon's Temple around 3000 years ago. They were said to be used to blow down the walls of Jericho. They are still used on certain religious days.[9]
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+ The Salpinx was a straight trumpet 62 inches (1,600 mm) long, made of bone or bronze. Salpinx contests were a part of the original Olympic Games.[9]
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+
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+ The Moche people of ancient Peru depicted trumpets in their art going back to AD 300.[10] The earliest trumpets were signaling instruments used for military or religious purposes, rather than music in the modern sense;[11] and the modern bugle continues this signaling tradition.
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+
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+ Improvements to instrument design and metal making in the late Middle Ages and Renaissance led to an increased usefulness of the trumpet as a musical instrument. The natural trumpets of this era consisted of a single coiled tube without valves and therefore could only produce the notes of a single overtone series. Changing keys required the player to change crooks of the instrument.[9] The development of the upper, "clarino" register by specialist trumpeters—notably Cesare Bendinelli—would lend itself well to the Baroque era, also known as the "Golden Age of the natural trumpet." During this period, a vast body of music was written for virtuoso trumpeters. The art was revived in the mid-20th century and natural trumpet playing is again a thriving art around the world. Many modern players in Germany and the UK who perform Baroque music use a version of the natural trumpet fitted with three or four vent holes to aid in correcting out-of-tune notes in the harmonic series.[12]
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+ The melody-dominated homophony of the classical and romantic periods relegated the trumpet to a secondary role by most major composers owing to the limitations of the natural trumpet. Berlioz wrote in 1844:
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+
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+ Notwithstanding the real loftiness and distinguished nature of its quality of tone, there are few instruments that have been more degraded (than the trumpet). Down to Beethoven and Weber, every composer – not excepting Mozart – persisted in confining it to the unworthy function of filling up, or in causing it to sound two or three commonplace rhythmical formulae.[13]
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+
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+ The trumpet is constructed of brass tubing bent twice into a rounded oblong shape.[14] As with all brass instruments, sound is produced by blowing air through closed lips, producing a "buzzing" sound into the mouthpiece and starting a standing wave vibration in the air column inside the trumpet. The player can select the pitch from a range of overtones or harmonics by changing the lip aperture and tension (known as the embouchure).
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+
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+ The mouthpiece has a circular rim, which provides a comfortable environment for the lips' vibration. Directly behind the rim is the cup, which channels the air into a much smaller opening (the back bore or shank) that tapers out slightly to match the diameter of the trumpet's lead pipe. The dimensions of these parts of the mouthpiece affect the timbre or quality of sound, the ease of playability, and player comfort. Generally, the wider and deeper the cup, the darker the sound and timbre.
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+
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+ Modern trumpets have three (or, infrequently, four) piston valves, each of which increases the length of tubing when engaged, thereby lowering the pitch. The first valve lowers the instrument's pitch by a whole step (two semitones), the second valve by a half step (one semitone), and the third valve by one and a half steps (three semitones). When a fourth valve is present, as with some piccolo trumpets, it usually lowers the pitch a perfect fourth (five semitones). Used singly and in combination these valves make the instrument fully chromatic, i.e., able to play all twelve pitches of classical music. For more information about the different types of valves, see Brass instrument valves.
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+
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+ The pitch of the trumpet can be raised or lowered by the use of the tuning slide. Pulling the slide out lowers the pitch; pushing the slide in raises it. To overcome the problems of intonation and reduce the use of the slide, Renold Schilke designed the tuning-bell trumpet. Removing the usual brace between the bell and a valve body allows the use of a sliding bell; the player may then tune the horn with the bell while leaving the slide pushed in, or nearly so, thereby improving intonation and overall response.[15]
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+
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+ A trumpet becomes a closed tube when the player presses it to the lips; therefore, the instrument only naturally produces every other overtone of the harmonic series. The shape of the bell makes the missing overtones audible.[16] Most notes in the series are slightly out of tune and modern trumpets have slide mechanisms for the first and third valves with which the player can compensate by throwing (extending) or retracting one or both slides, using the left thumb and ring finger for the first and third valve slides respectively.
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+
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+ The most common type is the B♭ trumpet, but A, C, D, E♭, E, low F, and G trumpets are also available. The C trumpet is most common in American orchestral playing, where it is used alongside the B♭ trumpet. Orchestral trumpet players are adept at transposing music at sight, frequently playing music written for the A, B♭, D, E♭, E, or F trumpet on the C trumpet or B♭ trumpet.
37
+
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+ The smallest trumpets are referred to as piccolo trumpets. The most common of these are built to play in both B♭ and A, with separate leadpipes for each key. The tubing in the B♭ piccolo trumpet is one-half the length of that in a standard B♭ trumpet. Piccolo trumpets in G, F and C are also manufactured, but are less common. Many players use a smaller mouthpiece on the piccolo trumpet, which requires a different sound production technique from the B♭ trumpet and can limit endurance. Almost all piccolo trumpets have four valves instead of the usual three — the fourth valve lowers the pitch, usually by a fourth, to assist in the playing of lower notes and to create alternate fingerings that facilitate certain trills. Maurice André, Håkan Hardenberger, David Mason, and Wynton Marsalis are some well-known trumpet players known for their additional virtuosity on the piccolo trumpet.
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+
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+ Trumpets pitched in the key of low G are also called sopranos, or soprano bugles, after their adaptation from military bugles. Traditionally used in drum and bugle corps, sopranos have featured both rotary valves and piston valves.
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+ The bass trumpet is usually played by a trombone player, being at the same pitch.[4] Bass trumpet is played with a shallower trombone mouthpiece, and music for it is written in treble clef. The most common keys for bass trumpets are C and B♭. Both C and B♭ bass trumpets are transposing instruments sounding an octave (C) or a major ninth (B♭) lower than written.
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+
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+ The historical slide trumpet was probably first developed in the late 14th century for use in alta cappella wind bands. Deriving from early straight trumpets, the Renaissance slide trumpet was essentially a natural trumpet with a sliding leadpipe. This single slide was rather awkward, as the entire corpus of the instrument moved, and the range of the slide was probably no more than a major third. Originals were probably pitched in D, to fit with shawms in D and G, probably at a typical pitch standard near A=466 Hz. As no known instruments from this period survive, the details—and even the existence—of a Renaissance slide trumpet is a matter of conjecture and debate among scholars.[17]
45
+
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+ Some slide trumpet designs saw use in England in the 18th century.[18]
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+
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+ The pocket trumpet is a compact B♭ trumpet. The bell is usually smaller than a standard trumpet and the tubing is more tightly wound to reduce the instrument size without reducing the total tube length. Its design is not standardized, and the quality of various models varies greatly. It can have a tone quality and projection unique in the trumpet world: a warm sound and a voice-like articulation. Since many pocket trumpet models suffer from poor design as well as cheap and imprecise manufacturing, the intonation, tone color and dynamic range of such instruments are severely hindered. Professional-standard instruments are, however, available. While they are not a substitute for the full-sized instrument, they can be useful in certain contexts. The jazz musician Don Cherry was renowned for his playing of the pocket instrument.
49
+
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+ The herald trumpet has an elongated bell extending far in front of the player, allowing a standard length of tubing from which a flag may be hung; the instrument is mostly used for ceremonial events such as parades and fanfares.
51
+
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+ Monette designed the flumpet in 1989 for jazz musician Art Farmer. It is a hybrid instrument with elements of trumpet and flugelhorn, sharing the three piston valve design and with a pitch of B♭.[19]
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+
54
+ There are also rotary-valve, or German, trumpets (which are commonly used in professional German and Austrian orchestras) as well as alto and Baroque trumpets.
55
+
56
+ Another variant of the standard trumpet is the Vienna valve trumpet. Primarily used in Viennese brass ensembles and orchestras such as the Vienna Philharmonic and Mnozil Brass.
57
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+ The trumpet is often confused with its close relative the cornet, which has a more conical tubing shape compared to the trumpet's more cylindrical tube. This, along with additional bends in the cornet's tubing, gives the cornet a slightly mellower tone, but the instruments are otherwise nearly identical. They have the same length of tubing and, therefore, the same pitch, so music written for cornet and trumpet is interchangeable. Another relative, the flugelhorn, has tubing that is even more conical than that of the cornet, and an even richer tone. It is sometimes augmented with a fourth valve to improve the intonation of some lower notes.
59
+
60
+ On any modern trumpet, cornet, or flugelhorn, pressing the valves indicated by the numbers below produces the written notes shown. "Open" means all valves up, "1" means first valve, "1–2" means first and second valve simultaneously, and so on. The sounding pitch depends on the transposition of the instrument. Engaging the fourth valve, if present, usually drops any of these pitches by a perfect fourth as well. Within each overtone series, the different pitches are attained by changing the embouchure. Standard fingerings above high C are the same as for the notes an octave below (C♯ is 1–2, D is 1, etc.)[20]
61
+
62
+ Each overtone series on the trumpet begins with the first overtone—the fundamental of each overtone series cannot be produced except as a pedal tone. Notes in parentheses are the sixth overtone, representing a pitch with a frequency of seven times that of the fundamental; while this pitch is close to the note shown, it is slightly flat relative to equal temperament, and use of those fingerings is generally avoided.
63
+
64
+ The fingering schema arises from the length of each valve's tubing (a longer tube produces a lower pitch). Valve "1" increases the tubing length enough to lower the pitch by one whole step, valve "2" by one half step, and valve "3" by one and a half steps. This scheme and the nature of the overtone series create the possibility of alternate fingerings for certain notes. For example, third-space "C" can be produced with no valves engaged (standard fingering) or with valves 2–3. Also, any note produced with 1–2 as its standard fingering can also be produced with valve 3 – each drops the pitch by ​1 1⁄2 steps. Alternate fingerings may be used to improve facility in certain passages, or to aid in intonation. Extending the third valve slide when using the fingerings 1–3 or 1-2-3 further lowers the pitch slightly to improve intonation.
65
+
66
+ If you take the partials of the harmonic series that a modern Bb trumpet can play for each combination of valves pressed, some notes are in tune with 12-tone equal temperament and some are not. The following tables show all notes from all partials from all valve combinations in order up to high C (C6). This often gives multiple fingerings for the same note and some notes that aren't usually useful.
67
+
68
+ Various types of mutes can be used to alter the sound of the instrument when placed in or over the bell. While most types of mutes do decrease the volume the instrument produces, as the name implies, the sound modification is typically the primary reason for their use. Types of mutes most commonly used to alter the sound of the instrument are: Straight Mutes, Harmon Mutes (aka "Wah-Wah" Mutes), Plunger Mutes, Bucket Mutes, and Cup Mutes. A description of their construction and sound quality are below:
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+
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+ Straight Mute: Constructed of either aluminum, which produces a bright piercing sound, or stone lined with cardboard, which produces a stuffy sound.
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+
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+ Harmon Mute: Constructed of aluminum and consists of two parts called the "stem" and the "body". The stem can be extended or removed to produce different timbres of sound. This mute is also called the "Wah-Wah" mute due to its distinctive sound created by the player placing their hand over the stem opening and waving it back and forth.
73
+
74
+ Plunger Mute: Most often made of a rubber bathroom plunger without the stick. This is used to manipulate sound by the player holding it over the bell with their left hand.
75
+
76
+ Bucket Mute: Constructed from cardboard and cloth, this mute is clipped to the end of the bell and used to muffle the sound almost completely.
77
+
78
+ Cup Mute: Also constructed of cardboard, this mute is shaped exactly like a straight mute but includes a cup at the end. In many models the cup is adjustable much like the stem on the harmon mute and produces a softer more muffled sound than a traditional straight mute.
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+ The standard trumpet range extends from the written F♯ immediately below Middle C up to about three octaves higher (F♯3 – F♯6). Traditional trumpet repertoire rarely calls for notes beyond this range, and the fingering tables of most method books peak at the high C, two octaves above middle C.[contradictory] Several trumpeters have achieved fame for their proficiency in the extreme high register, among them Maynard Ferguson, Cat Anderson, Dizzy Gillespie, Doc Severinsen, and more recently Wayne Bergeron, Thomas Gansch, James Morrison, Jon Faddis and Arturo Sandoval. It is also possible to produce pedal tones below the low F♯, which is a device occasionally employed in the contemporary repertoire for the instrument.
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+
82
+ Contemporary music for the trumpet makes wide uses of extended trumpet techniques.
83
+
84
+ Flutter tonguing: The trumpeter rolls the tip of the tongue to produce a 'growling like' tone. It is achieved as if one were rolling an R in the Spanish language. This technique is widely employed by composers like Berio and Stockhausen.
85
+
86
+ Growling: Simultaneously playing tone while using the back of the tongue to vibrate the uvula creating a distinct sound. Most trumpet players will use a plunger with this technique to achieve a particular sound heard in a lot of Chicago Jazz of the 1950s.
87
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+ Double tonguing: The player articulates using the syllables ta-ka ta-ka ta-ka
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+ Triple tonguing: The same as double tonguing, but with the syllables ta-ta-ka ta-ta-ka ta-ta-ka or ta-ka-ta ta-ka-ta.
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+ Doodle tongue: The trumpeter tongues as if saying the word doodle. This is a very faint tonguing similar in sound to a valve tremolo.
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+ Glissando: Trumpeters can slide between notes by depressing the valves halfway and changing the lip tension. Modern repertoire makes extensive use of this technique.
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+ Vibrato: It is often regulated in contemporary repertoire through specific notation. Composers can call for everything from fast, slow or no vibrato to actual rhythmic patterns played with vibrato.
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+ Pedal tone: Composers have written for two-and-a-half octaves below the low F♯, which is at the bottom of the standard range. Extreme low pedals are produced by slipping the lower lip out of the mouthpiece. Claude Gordon assigned pedals as part of his trumpet practice routines, that were a systematic expansion on his lessons with Herbert L. Clarke. The technique was pioneered by Bohumir Kryl.[21]
99
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+ Microtones: Composers such as Scelsi and Stockhausen have made wide use of the trumpet's ability to play microtonally. Some instruments feature a fourth valve that provides a quarter-tone step between each note. The jazz musician Ibrahim Maalouf uses such a trumpet, invented by his father to make it possible to play Arab maqams.
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+ Valve tremolo: Many notes on the trumpet can be played in several different valve combinations. By alternating between valve combinations on the same note, a tremolo effect can be created. Berio makes extended use of this technique in his Sequenza X.
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+ Noises: By hissing, clicking, or breathing through the instrument, the trumpet can be made to resonate in ways that do not sound at all like a trumpet. Noises may require amplification.
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+ Preparation: Composers have called for trumpeters to play under water, or with certain slides removed. It is increasingly common for composers to specify all sorts of preparations for trumpet. Extreme preparations involve alternate constructions, such as double bells and extra valves.
107
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+ Split tone: Trumpeters can produce more than one tone simultaneously by vibrating the two lips at different speeds. The interval produced is usually an octave or a fifth.
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+ Lip-trill or shake: Also known as "lip-slurs". By rapidly varying air speed, but not changing the depressed valves, the pitch varies quickly between adjacent harmonic partials. Shakes and lip-trills can vary in speed and the distance between the partials can be as large or small as the musicians' desires. Traditionally, however, lip-trills and shakes are usually the next partial up from the written note.
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+ Multi-phonics: Playing a note and "humming" a different note simultaneously. For example, sustaining a middle C and humming a major 3rd "E" at the same time.
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+ Circular breathing: A technique wind players use to produce uninterrupted tone, without pauses for breaths. The player puffs up the cheeks, storing air, then breathes in rapidly through the nose while using the cheeks to continue pushing air outwards.
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+ One trumpet method is Jean-Baptiste Arban's Complete Conservatory Method for Trumpet (Cornet).[22] Other well-known method books include Technical Studies by Herbert L. Clarke,[23] Grand Method by Louis Saint-Jacome, Daily Drills and Technical Studies by Max Schlossberg, and methods by Ernest S. Williams, Claude Gordon, Charles Colin, James Stamp, and Louis Davidson.[24] A common method book for beginners is the Walter Beeler's Method for the Cornet, and there have been several instruction books written by virtuoso Allen Vizzutti.[25] Merri Franquin wrote a Complete Method for Modern Trumpet,[26] which fell into obscurity for much of the twentieth century until public endorsements by Maurice André revived interest in this work.[27]
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+ In early jazz, Louis Armstrong was well known for his virtuosity and his improvisations on the Hot Five and Hot Seven recordings, and his switch from cornet to trumpet is often cited as heralding the trumpet's dominance over the cornet in jazz.[4][28] Dizzy Gillespie was a gifted improviser with an extremely high (but musical) range, building on the style of Roy Eldridge but adding new layers of harmonic complexity. Gillespie had an enormous impact on virtually every subsequent trumpeter, both by the example of his playing and as a mentor to younger musicians. Miles Davis is widely considered one of the most influential musicians of the 20th century—his style was distinctive and widely imitated. Davis' phrasing and sense of space in his solos have been models for generations of jazz musicians.[29] Cat Anderson was a trumpet player who was known for the ability to play extremely high with an even more extreme volume, who played with Duke Ellington's Big Band. Maynard Ferguson came to prominence playing in Stan Kenton's orchestra, before forming his own band in 1957. He was noted for being able to play accurately in a remarkably high register.[30]
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+
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+ Anton Weidinger developed in the 1790s the first successful keyed trumpet, capable of playing all the chromatic notes in its range. Joseph Haydn's Trumpet Concerto was written for him in 1796 and startled contemporary audiences by its novelty,[31] a fact shown off by some stepwise melodies played low in the instrument's range.
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+
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+ The Last Judgment (workshop of Hieronymus Bosch) [nl], c.1500-1510
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+
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+ Trumpet-Player in front of a Banquet, Gerrit Dou, c.1660-1665
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+
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+ Illustration for The Trumpeter Taken Prisoner from Baby's Own Aesop, a children's edition of Aesop's fables
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+ Louis Armstrong statue in Algiers, New Orleans
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+ Miles Davis statue in Kielce, Poland
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1
+ Tropical climate is one of the five major climate groups in the Köppen climate classification. Tropical climates are characterized by monthly average temperatures of 18 ℃ (64.4 ℉) or higher year-round and feature hot temperatures. Annual precipitation is often abundant in tropical climates, and show a seasonal rhythm to varying degrees. Tropical climates are generally located within 20 to 25 degrees of the equator. Sunlight is intense.
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+ There are three basic types of tropical climates within the tropical climate group: tropical rainforest climate (Af), tropical monsoon climate (Am) and tropical wet and dry or savanna climate (Aw or As), which are classified and distinguished by the annual precipitation and the precipitation level of the driest month in those regions.
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+ Köppen climate classification is one of the most widely used climate classification systems. It defines a tropical climate as a region where the temperature of the coldest month is greater than or equal to 18 ℃ (64.4 ℉) and classifies them as an A-group (tropical climate group).[1] A-group regions are usually found around the Equator, including Central America, northern part of South America, central part of Africa, southern part of Asia and the Pacific Ocean islands.[2]
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+ In Group A, there are three types of climate: tropical rainforest climate (Af), tropical monsoon climate (Am) and tropical wet and dry or savanna climate (Aw or As). All of the three climates are classified by their Pdry (short for precipitation of the driest month). Tropical rainforest climate’s Pdry should be greater or equal than 60 mm (2.4 inches). Tropical monsoon climate’s Pdry should be in the range from
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+ 100
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+ −
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+ {\displaystyle 100-{\tfrac {mean\ annual\ precipitation}{25}}}
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+ to 60 mm. Tropical wet and dry or savanna climate’s Pdry should be less than
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+ −
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+ 25
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+ {\displaystyle 100-{\tfrac {mean\ annual\ precipitation}{25}}}
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+ .[1]
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+ The tropics have the characteristics of small temperature changes and long summers. Due to the high temperature and abundant rainfall, some plants can grow throughout the year. High temperature and humidity is the most suitable environment for epiphytes to grow. Plants of all sizes can vegetate under tropical climates. Vegetations grow in layers: shrubs under tall trees, and bushes under shrubs. Almost every inch of space is being well used. Tropical plants are rich in resources, including coffee, cocoa and oil palm.[3][4] Listed below are types of vegetation unique to each of the three climates that make up the tropical climate biome.
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+ Tropical rainforest vegetation including: Bengal bamboo, bougainvillea, curare, coconut tree, Durian
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+ Tropical monsoon vegetation including: teak, deodar, rosewood, sandalwood and bamboo
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+ Tropical wet and dry or savanna vegetation including: acacia senegal, elephant grass, jarrah tree, gum tree eucalyptus, whistling thorn
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+
105
+ The Köppen classification identifies tropical rainforest climate (Zone Af: f = feucht, German for moist) as having north and south latitudinal ranges of 5-10 degrees from the equator.[5][6] Tropical rainforest climates have high temperatures: the yearly average temperature is between 70 °F to 85 °F (21 °C to 30 °C).[7][8] The precipitation can reach over 100 inches a year.[7][8] The seasons are evenly distributed throughout the year, and there is almost no drought period.[6] Regions affected by tropical rainforest climate mainly include the upper Amazon basin of South America, the Northern Zaire (Congo) basin of Africa, and the islands of the East Indies.[6]
106
+
107
+ The tropical rainforest climate differs from other subtypes of tropical climates as it has more kinds of trees.[8] The large number of trees contribute to the humidity of the climate because of the transpiration, which is the process of water lost from the surface of living plants to the atmosphere. The warmth and abundant precipitation contributes to the diversity and characteristics of vegetations under the tropical rainforest climate.[7] The vegetations develop a vertical stratification and various growth forms to receive enough sunlight, which is unusual under other types of climate.[7]
108
+
109
+ The Köppen classification tool identifies tropical monsoon climate as having small annual temperature ranges, high temperatures, and plentiful precipitation. This climate also has a short dry season which occurs in the winter.[9] The tropical monsoon climate is usually found within countries in the south and southeast Asia region between the latitude of 10 degrees north and the Tropic of Cancer. These regions include India, Philippines, the northern coast of Australia and Hainan Island of China. The annual temperature of regions under tropical monsoon climate is stable.
110
+
111
+ The tropical monsoon climate has the following main characteristic. The average annual temperature is 27.05 °C (80.69 °F) and has an annual temperature range of 3.6°C (2°F).[10] Distinction between wet and drought seasons, the tropical monsoon climate is different from other tropical climates because of its uneven precipitation throughout the year. The precipitation is heavy in the summer, and a short-drought season occurs in the winter. This climate has an annual total precipitation of 3409.2mm, and a 3115.9mm summer precipitation and 293.3mm winter precipitation.[10]
112
+
113
+ There are three main seasons of tropical monsoon climate: the cool dry season is from October to February, the hot dry season is from March to mid-June and the rainy season is from mid-June to September.[11]
114
+
115
+ The tropical monsoon forest mainly consists of three layered structures. The first layer is the ground layer which is a very dense layer of shrubs. The second layer is the understory layer with trees about 15 meters tall. The top layer is called the canopy tree which has trees from 25 to 30 meters tall and those trees grow closely.[citation needed]
116
+
117
+ The regions under this climate are mainly located near the equator between 10° north-south latitude and north and south of the tropic of cancer. Typical regions include central Africa, most of South America, as well as northern and eastern Australia.[12] The temperature range of savanna climate is between 20 °C to 30 °C (68 °F - 86 °F). In summer, the temperature is between 25 °C - 30 °C, while in winter the temperature is between 20 °C - 30 °C.[13] The annual precipitation is between 700 to 1000 mm. The driest months are from November to March and they have less than 60 mm of rainfall.
118
+
119
+ Regions under the savanna climate usually have lands covered with flat grassland vegetation. Those grassland biomes cover almost 20% of the earth’s surface.[14] The grassland vegetation types include Rhodes grass, red oats grass, star grass and lemongrass.[15]
en/5806.html.txt ADDED
@@ -0,0 +1,37 @@
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
1
+ Coordinates: 23°26′12.5″N 0°0′0″W / 23.436806°N -0.00000°E / 23.436806; -0.00000 (Prime Meridian)
2
+
3
+ The Tropic of Cancer, which is also referred to as the Northern Tropic, is the most northerly circle of latitude on Earth at which the Sun can be directly overhead. This occurs on the June solstice, when the Northern Hemisphere is tilted toward the Sun to its maximum extent.[1] It also reaches 90 degrees below the horizon at solar midnight on the December Solstice. Using a continuously updated formula, the circle is currently 23°26′11.8″ (or 23.43662°) north of the Equator.
4
+
5
+ Its Southern Hemisphere counterpart, marking the most southerly position at which the Sun can be directly overhead, is the Tropic of Capricorn. These tropics are two of the five major circles of latitude that mark maps of Earth, the others being the Arctic and Antarctic Circles and the Equator. The positions of these two circles of latitude (relative to the Equator) are dictated by the tilt of Earth's axis of rotation relative to the plane of its orbit, and since the tilt changes, the location of these two circles also changes.
6
+
7
+ When this line of latitude was named in the last centuries BC, the Sun was in the constellation Cancer (Latin for crab) at the June solstice, the time each year that the Sun reaches its zenith at this latitude. Due to the precession of the equinoxes, this is no longer the case; today the Sun is in Taurus at the June solstice. The word "tropic" itself comes from the Greek "trope (τροπή)", meaning turn (change of direction, or circumstances), inclination, referring to the fact that the Sun appears to "turn back" at the solstices.
8
+
9
+ The Tropic of Cancer's position is not fixed, but constantly changes because of a slight wobble in the Earth's longitudinal alignment relative to the ecliptic, the plane in which the Earth orbits around the Sun. Earth's axial tilt varies over a 41,000-year period from 22.1 to 24.5 degrees, and as of 2000[update] is about 23.4 degrees, which will continue to remain valid for about a millennium. This wobble means that the Tropic of Cancer is currently drifting southward at a rate of almost half an arcsecond (0.468″) of latitude, or 15 metres, per year. The circle's position was at exactly 23° 27′N in 1917 and will be at 23° 26'N in 2045.[2]
10
+
11
+ See axial tilt and circles of latitude for additional details.
12
+
13
+ North of the tropic are the subtropics and the North Temperate Zone. The equivalent line of latitude south of the Equator is called the Tropic of Capricorn, and the region between the two, centered on the Equator, is the tropics.
14
+
15
+ There are approximately 13 hours, 35 minutes of daylight during the summer solstice. During the winter solstice, there are 10 hours, 41 minutes of daylight.
16
+
17
+ Using 23°26'N for the Tropic of Cancer, the tropic passes through the following countries and territories starting at the prime meridian and heading eastward:
18
+
19
+ Excluding cooler highland regions in China, the climate at the Tropic of Cancer is generally hot and dry except for easterly coastal areas where orographic rainfall can be very heavy, in some places reaching 4 metres (160 in) annually. Most regions on the Tropic of Cancer experience two distinct seasons: an extremely hot summer with temperatures often reaching 45 °C (113 °F) and a warm winter with maxima around 22 °C (72 °F). Much land on or near the Tropic of Cancer is part of the Sahara Desert, while to the east the climate is torrid monsoonal with a short wet season from June to September and very little rainfall for the rest of the year.
20
+
21
+ The highest mountain on or adjacent to the Tropic of Cancer is Yu Shan in Taiwan; though it had glaciers descending as low as 2,800 metres (9,190 ft) during the Last Glacial Maximum, none survive and at present no glaciers exist within 470 kilometres (290 mi) of the Tropic of Cancer; the nearest currently surviving are the Minyong and Baishui in the Himalayas to the north and on Iztaccíhuatl to the south.
22
+
23
+ According to the rules of the Fédération Aéronautique Internationale, for a flight to compete for a round-the-world speed record, it must cover a distance no less than the length of the Tropic of Cancer, cross all meridians, and end on the same airfield where it started.
24
+
25
+ Length of the Tropic of Cancer at 23°26′11.8″N (4 July 2020 to 19 September 2020) is 36,788 kilometres (22,859 mi).[3]
26
+
27
+ For an ordinary circumnavigation the rules are somewhat relaxed and the distance is set to a rounded value of at least 36,770 kilometres (22,850 mi).
28
+
29
+ Road sign south of Dakhla, Western Sahara marking the Tropic of Cancer. The sign was placed by Budapest-Bamako rally participants; thus, the inscription is in English and Hungarian.
30
+
31
+ Sign marking the Tropic of Cancer a few miles from Rann of Kutch, Gujarat, India
32
+
33
+ Sign marking the Tropic of Cancer in Madhya Pradesh, India
34
+
35
+ Sign marking the Tropic of Cancer on National Highway 34 in Nadia District, West Bengal, India
36
+
37
+ Ruisui Tropic of Cancer Marker in Ruisui Township, Hualien County, Taiwan
en/5807.html.txt ADDED
@@ -0,0 +1,26 @@
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
1
+ The tropics is the region of the Earth surrounding the Equator. They are delimited in latitude by the Tropic of Cancer in the Northern Hemisphere at 23°26′11.8″ (or 23.43662°) N and the Tropic of Capricorn in
2
+ the Southern Hemisphere at 23°26′11.8″ (or 23.43662°) S; these latitudes correspond to the axial tilt of the Earth. The tropics are also referred to as the tropical zone and the torrid zone (see geographical zone). The tropics include all the areas on the Earth where the Sun contacts a point directly overhead at least once during the solar year (which is a subsolar point) - thus the latitude of the tropics is roughly equal to the angle of the Earth's axial tilt.
3
+
4
+ In terms of climate, the tropics receive sunlight that is more direct than the rest of Earth and are generally hotter and wetter. The word "tropical" sometimes refers to this sort of climate rather than to the geographical zone. The tropical zone includes deserts and snow-capped mountains, which are not tropical in the climatic sense. The tropics are distinguished from the other climatic and biomatic regions of Earth, which are the middle latitudes and the polar regions on either side of the equatorial zone.
5
+
6
+ The tropics constitute 40% of the Earth's surface area[1] and contain 36% of the Earth's landmass.[2] As of 2014[update], the region was home to 40% of the world's population, and this figure was then projected to reach 50% by 2050.[3]
7
+
8
+ The word "tropic" comes from Ancient Greek τροπή (tropē), meaning "to turn" or "change direction", as the Sun appears to cease its southerly course when it reaches this latitude and begins moving back to the north.
9
+
10
+ "Tropical" is sometimes used in a general sense for a tropical climate to mean warm to hot and moist year-round, often with the sense of lush vegetation.
11
+
12
+ Many tropical areas have a dry and wet season. The wet season, rainy season or green season is the time of year, ranging from one or more months, when most of the average annual rainfall in a region falls.[4] Areas with wet seasons are disseminated across portions of the tropics and subtropics.[5] Under the Köppen climate classification, for tropical climates, a wet-season month is defined as a month where average precipitation is 60 millimetres (2.4 in) or more.[6] Tropical rainforests technically do not have dry or wet seasons, since their rainfall is equally distributed through the year.[7] Some areas with pronounced rainy seasons see a break in rainfall during mid-season when the intertropical convergence zone or monsoon trough moves poleward of their location during the middle of the warm season;[8] typical vegetation in these areas ranges from moist seasonal tropical forests to savannahs.
13
+
14
+ When the wet season occurs during the warm season, or summer, precipitation falls mainly during the late afternoon and early evening hours. The wet season is a time when air quality improves, freshwater quality improves and vegetation grows significantly, leading to crop yields late in the season. Floods cause rivers to overflow their banks, and some animals to retreat to higher ground. Soil nutrients diminish and erosion increases. The incidence of malaria increases in areas where the rainy season coincides with high temperatures. Animals have adaptation and survival strategies for the wetter regime. The previous dry season leads to food shortages into the wet season, as the crops have yet to mature.
15
+
16
+ However, regions within the tropics may well not have a tropical climate. Under the Köppen climate classification, much of the area within the geographical tropics is classed not as "tropical" but as "dry" (arid or semi-arid), including the Sahara Desert, the Atacama Desert and Australian Outback. Also, there are alpine tundra and snow-capped peaks, including Mauna Kea, Mount Kilimanjaro, and the Andes as far south as the northernmost parts of Chile and Perú.
17
+
18
+ Tropical plants and animals are those species native to the tropics. Tropical ecosystems may consist of tropical rainforests, seasonal tropical forests, dry (often deciduous) forests, spiny forests, desert and other habitat types. There are often significant areas of biodiversity, and species endemism present, particularly in rainforests and seasonal forests. Some examples of important biodiversity and high endemism ecosystems are El Yunque National Forest in Puerto Rico, Costa Rican and Nicaraguan rainforests, Amazon Rainforest territories of several South American countries, Madagascar dry deciduous forests, the Waterberg Biosphere of South Africa, and eastern Madagascar rainforests. Often the soils of tropical forests are low in nutrient content, making them quite vulnerable to slash-and-burn deforestation techniques, which are sometimes an element of shifting cultivation agricultural systems.
19
+
20
+ In biogeography, the tropics are divided into Paleotropics (Africa, Asia and Australia) and Neotropics (Caribbean, Central America, and South America). Together, they are sometimes referred to as the Pantropic. The system of biogeographic realms differs somewhat; the Neotropical realm includes both the Neotropics and temperate South America, and the Paleotropics correspond to the Afrotropical, Indomalayan, Oceanian, and tropical Australasian realms.
21
+
22
+ Tropicality refers to the image that people outside the tropics have of the region, ranging from critical to verging on fetishism. The idea of tropicality gained renewed interest in geographical discourse when French geographer Pierre Gourou published Les Pays Tropicaux (The Tropical World in English), in the late 1940s.[9]
23
+
24
+ Tropicality encompassed two images. One, is that the tropics represent a 'Garden of Eden', a heaven on Earth, a land of rich biodiversity - aka a tropical paradise.[10] The alternative is that the tropics consist of wild, unconquerable nature. The latter view was often discussed in old Western literature more so than the first.[10] Evidence suggests over time that the view of the tropics as such in popular literature has been supplanted by more well-rounded and sophisticated interpretations.[11]
25
+
26
+ Western scholars tried to theorize reasons about why tropical areas were relatively more inhospitable to human civilisations then those existing in colder regions of the Northern Hemisphere. A popular explanation focused on the differences in climate. Tropical jungles and rainforests have much more humid and hotter weather than colder and drier temperaments of the Northern Hemisphere. This theme led to some scholars to suggest that humid hot climates correlate to human populations lacking control over nature e.g. ' the wild Amazonian rainforests'.[12]
en/5808.html.txt ADDED
@@ -0,0 +1,37 @@
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
1
+ Coordinates: 23°26′12.5″N 0°0′0″W / 23.436806°N -0.00000°E / 23.436806; -0.00000 (Prime Meridian)
2
+
3
+ The Tropic of Cancer, which is also referred to as the Northern Tropic, is the most northerly circle of latitude on Earth at which the Sun can be directly overhead. This occurs on the June solstice, when the Northern Hemisphere is tilted toward the Sun to its maximum extent.[1] It also reaches 90 degrees below the horizon at solar midnight on the December Solstice. Using a continuously updated formula, the circle is currently 23°26′11.8″ (or 23.43662°) north of the Equator.
4
+
5
+ Its Southern Hemisphere counterpart, marking the most southerly position at which the Sun can be directly overhead, is the Tropic of Capricorn. These tropics are two of the five major circles of latitude that mark maps of Earth, the others being the Arctic and Antarctic Circles and the Equator. The positions of these two circles of latitude (relative to the Equator) are dictated by the tilt of Earth's axis of rotation relative to the plane of its orbit, and since the tilt changes, the location of these two circles also changes.
6
+
7
+ When this line of latitude was named in the last centuries BC, the Sun was in the constellation Cancer (Latin for crab) at the June solstice, the time each year that the Sun reaches its zenith at this latitude. Due to the precession of the equinoxes, this is no longer the case; today the Sun is in Taurus at the June solstice. The word "tropic" itself comes from the Greek "trope (τροπή)", meaning turn (change of direction, or circumstances), inclination, referring to the fact that the Sun appears to "turn back" at the solstices.
8
+
9
+ The Tropic of Cancer's position is not fixed, but constantly changes because of a slight wobble in the Earth's longitudinal alignment relative to the ecliptic, the plane in which the Earth orbits around the Sun. Earth's axial tilt varies over a 41,000-year period from 22.1 to 24.5 degrees, and as of 2000[update] is about 23.4 degrees, which will continue to remain valid for about a millennium. This wobble means that the Tropic of Cancer is currently drifting southward at a rate of almost half an arcsecond (0.468″) of latitude, or 15 metres, per year. The circle's position was at exactly 23° 27′N in 1917 and will be at 23° 26'N in 2045.[2]
10
+
11
+ See axial tilt and circles of latitude for additional details.
12
+
13
+ North of the tropic are the subtropics and the North Temperate Zone. The equivalent line of latitude south of the Equator is called the Tropic of Capricorn, and the region between the two, centered on the Equator, is the tropics.
14
+
15
+ There are approximately 13 hours, 35 minutes of daylight during the summer solstice. During the winter solstice, there are 10 hours, 41 minutes of daylight.
16
+
17
+ Using 23°26'N for the Tropic of Cancer, the tropic passes through the following countries and territories starting at the prime meridian and heading eastward:
18
+
19
+ Excluding cooler highland regions in China, the climate at the Tropic of Cancer is generally hot and dry except for easterly coastal areas where orographic rainfall can be very heavy, in some places reaching 4 metres (160 in) annually. Most regions on the Tropic of Cancer experience two distinct seasons: an extremely hot summer with temperatures often reaching 45 °C (113 °F) and a warm winter with maxima around 22 °C (72 °F). Much land on or near the Tropic of Cancer is part of the Sahara Desert, while to the east the climate is torrid monsoonal with a short wet season from June to September and very little rainfall for the rest of the year.
20
+
21
+ The highest mountain on or adjacent to the Tropic of Cancer is Yu Shan in Taiwan; though it had glaciers descending as low as 2,800 metres (9,190 ft) during the Last Glacial Maximum, none survive and at present no glaciers exist within 470 kilometres (290 mi) of the Tropic of Cancer; the nearest currently surviving are the Minyong and Baishui in the Himalayas to the north and on Iztaccíhuatl to the south.
22
+
23
+ According to the rules of the Fédération Aéronautique Internationale, for a flight to compete for a round-the-world speed record, it must cover a distance no less than the length of the Tropic of Cancer, cross all meridians, and end on the same airfield where it started.
24
+
25
+ Length of the Tropic of Cancer at 23°26′11.8″N (4 July 2020 to 19 September 2020) is 36,788 kilometres (22,859 mi).[3]
26
+
27
+ For an ordinary circumnavigation the rules are somewhat relaxed and the distance is set to a rounded value of at least 36,770 kilometres (22,850 mi).
28
+
29
+ Road sign south of Dakhla, Western Sahara marking the Tropic of Cancer. The sign was placed by Budapest-Bamako rally participants; thus, the inscription is in English and Hungarian.
30
+
31
+ Sign marking the Tropic of Cancer a few miles from Rann of Kutch, Gujarat, India
32
+
33
+ Sign marking the Tropic of Cancer in Madhya Pradesh, India
34
+
35
+ Sign marking the Tropic of Cancer on National Highway 34 in Nadia District, West Bengal, India
36
+
37
+ Ruisui Tropic of Cancer Marker in Ruisui Township, Hualien County, Taiwan
en/5809.html.txt ADDED
@@ -0,0 +1,37 @@
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
1
+ Coordinates: 23°26′12.5″N 0°0′0″W / 23.436806°N -0.00000°E / 23.436806; -0.00000 (Prime Meridian)
2
+
3
+ The Tropic of Cancer, which is also referred to as the Northern Tropic, is the most northerly circle of latitude on Earth at which the Sun can be directly overhead. This occurs on the June solstice, when the Northern Hemisphere is tilted toward the Sun to its maximum extent.[1] It also reaches 90 degrees below the horizon at solar midnight on the December Solstice. Using a continuously updated formula, the circle is currently 23°26′11.8″ (or 23.43662°) north of the Equator.
4
+
5
+ Its Southern Hemisphere counterpart, marking the most southerly position at which the Sun can be directly overhead, is the Tropic of Capricorn. These tropics are two of the five major circles of latitude that mark maps of Earth, the others being the Arctic and Antarctic Circles and the Equator. The positions of these two circles of latitude (relative to the Equator) are dictated by the tilt of Earth's axis of rotation relative to the plane of its orbit, and since the tilt changes, the location of these two circles also changes.
6
+
7
+ When this line of latitude was named in the last centuries BC, the Sun was in the constellation Cancer (Latin for crab) at the June solstice, the time each year that the Sun reaches its zenith at this latitude. Due to the precession of the equinoxes, this is no longer the case; today the Sun is in Taurus at the June solstice. The word "tropic" itself comes from the Greek "trope (τροπή)", meaning turn (change of direction, or circumstances), inclination, referring to the fact that the Sun appears to "turn back" at the solstices.
8
+
9
+ The Tropic of Cancer's position is not fixed, but constantly changes because of a slight wobble in the Earth's longitudinal alignment relative to the ecliptic, the plane in which the Earth orbits around the Sun. Earth's axial tilt varies over a 41,000-year period from 22.1 to 24.5 degrees, and as of 2000[update] is about 23.4 degrees, which will continue to remain valid for about a millennium. This wobble means that the Tropic of Cancer is currently drifting southward at a rate of almost half an arcsecond (0.468″) of latitude, or 15 metres, per year. The circle's position was at exactly 23° 27′N in 1917 and will be at 23° 26'N in 2045.[2]
10
+
11
+ See axial tilt and circles of latitude for additional details.
12
+
13
+ North of the tropic are the subtropics and the North Temperate Zone. The equivalent line of latitude south of the Equator is called the Tropic of Capricorn, and the region between the two, centered on the Equator, is the tropics.
14
+
15
+ There are approximately 13 hours, 35 minutes of daylight during the summer solstice. During the winter solstice, there are 10 hours, 41 minutes of daylight.
16
+
17
+ Using 23°26'N for the Tropic of Cancer, the tropic passes through the following countries and territories starting at the prime meridian and heading eastward:
18
+
19
+ Excluding cooler highland regions in China, the climate at the Tropic of Cancer is generally hot and dry except for easterly coastal areas where orographic rainfall can be very heavy, in some places reaching 4 metres (160 in) annually. Most regions on the Tropic of Cancer experience two distinct seasons: an extremely hot summer with temperatures often reaching 45 °C (113 °F) and a warm winter with maxima around 22 °C (72 °F). Much land on or near the Tropic of Cancer is part of the Sahara Desert, while to the east the climate is torrid monsoonal with a short wet season from June to September and very little rainfall for the rest of the year.
20
+
21
+ The highest mountain on or adjacent to the Tropic of Cancer is Yu Shan in Taiwan; though it had glaciers descending as low as 2,800 metres (9,190 ft) during the Last Glacial Maximum, none survive and at present no glaciers exist within 470 kilometres (290 mi) of the Tropic of Cancer; the nearest currently surviving are the Minyong and Baishui in the Himalayas to the north and on Iztaccíhuatl to the south.
22
+
23
+ According to the rules of the Fédération Aéronautique Internationale, for a flight to compete for a round-the-world speed record, it must cover a distance no less than the length of the Tropic of Cancer, cross all meridians, and end on the same airfield where it started.
24
+
25
+ Length of the Tropic of Cancer at 23°26′11.8″N (4 July 2020 to 19 September 2020) is 36,788 kilometres (22,859 mi).[3]
26
+
27
+ For an ordinary circumnavigation the rules are somewhat relaxed and the distance is set to a rounded value of at least 36,770 kilometres (22,850 mi).
28
+
29
+ Road sign south of Dakhla, Western Sahara marking the Tropic of Cancer. The sign was placed by Budapest-Bamako rally participants; thus, the inscription is in English and Hungarian.
30
+
31
+ Sign marking the Tropic of Cancer a few miles from Rann of Kutch, Gujarat, India
32
+
33
+ Sign marking the Tropic of Cancer in Madhya Pradesh, India
34
+
35
+ Sign marking the Tropic of Cancer on National Highway 34 in Nadia District, West Bengal, India
36
+
37
+ Ruisui Tropic of Cancer Marker in Ruisui Township, Hualien County, Taiwan
en/581.html.txt ADDED
@@ -0,0 +1,49 @@
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
1
+
2
+
3
+ The bass guitar, electric bass, or simply bass, is the lowest-pitched member of the guitar family. It is a plucked string instrument similar in appearance and construction to an electric or an acoustic guitar, but with a longer neck and scale length, and typically four to six strings or courses. Since the mid-1950s, the bass guitar has largely replaced the double bass in popular music.
4
+
5
+ The four-string bass is usually tuned the same as the double bass, which corresponds to pitches one octave lower than the four lowest-pitched strings of a guitar (E, A, D, and G). It is played primarily with the fingers or thumb, or by striking with a pick. The electric bass guitar has pickups and must be connected to an amplifier and speaker.
6
+
7
+ According to the New Grove Dictionary of Music and Musicians, an "Electric bass guitar [is] a Guitar, usually with four heavy strings tuned E1'–A1'–D2–G2."[1] It also defines bass as "Bass (iv). A contraction of Double bass or Electric bass guitar." According to some authors the proper term is "electric bass".[2][3] Common names for the instrument are "bass guitar", "electric bass guitar", and "electric bass"[4] and some authors claim that they are historically accurate.[5] As the electric alternative to a double bass (which is not a guitar), many manufacturers such as Fender list the instrument in the electric bass category rather than the guitar category.[6]
8
+
9
+ The bass guitar is a transposing instrument, as it is notated in bass clef an octave higher than it sounds, to reduce the need for ledger lines in music written for the instrument, and simplify reading.[7]
10
+
11
+ In the 1930s, musician and inventor Paul Tutmarc of Seattle, Washington, developed the first electric bass guitar in its modern form, a fretted instrument designed to be played horizontally. The 1935 sales catalog for Tutmarc's company Audiovox featured his "Model 736 Bass Fiddle", a solid-bodied electric bass guitar with four strings, a 30 1⁄2-inch (775-millimetre) scale length, and a single pickup.[8] Around 100 were made during this period.[9] Audiovox also sold their “Model 236” bass amplifier.[10]
12
+
13
+ In the 1950s, Leo Fender and George Fullerton developed the first mass-produced electric bass guitar.[11] The Fender Electric Instrument Manufacturing Company began producing the Precision Bass, or P-Bass, in October 1951. The design featured a simple uncontoured "slab" body design and a single coil pickup similar to that of a Telecaster. By 1957 the Precision more closely resembled the Fender Stratocaster with the body edges beveled for comfort, and the pickup was changed to a split coil design.[12]
14
+
15
+ The Fender Bass was a revolutionary instrument for gigging musicians. In comparison with the large, heavy upright bass, which had been the main bass instrument in popular music from the early 20th century to the 1940s, the bass guitar could be easily transported to shows. When amplified, the bass guitar was also less prone than acoustic basses to unwanted audio feedback.[13] The addition of frets enabled bassists to play in tune more easily than on fretless acoustic or electric upright basses, and allowed guitarists to more easily transition to the instrument.[14]
16
+
17
+ In 1953, Monk Montgomery became the first bassist to tour with the Fender bass, in Lionel Hampton's postwar big band.[15] Montgomery was also possibly the first to record with the electric bass, on July 2, 1953, with the Art Farmer Septet.[16] Roy Johnson (with Lionel Hampton), and Shifty Henry (with Louis Jordan and His Tympany Five), were other early Fender bass pioneers.[11] Bill Black, who played with Elvis Presley, switched from upright bass to the Fender Precision Bass around 1957.[17] The bass guitar was intended to appeal to guitarists as well as upright bass players, and many early pioneers of the instrument, such as Carol Kaye, Joe Osborn, and Paul McCartney were originally guitarists.[13]
18
+
19
+ Also in 1953, Gibson released the first short-scale violin-shaped electric bass, with an extendable end pin so a bassist could play it upright or horizontally. Gibson renamed the bass the EB-1 in 1958.[citation needed] In 1958, Gibson released the maple arched-top EB-2 described in the Gibson catalog as a "hollow-body electric bass that features a Bass/Baritone pushbutton for two different tonal characteristics".[citation needed] In 1959, these were followed by the more conventional-looking EB-0 Bass.[citation needed] The EB-0 was very similar to a Gibson SG in appearance (although the earliest examples have a slab-sided body shape closer to that of the double-cutaway Les Paul Special). Whereas Fender basses had pickups mounted in positions in between the base of the neck and the top of the bridge, many of Gibson's early basses featured one humbucking pickup mounted directly against the neck pocket.[citation needed] The Fender and Gibson versions used bolt-on and glued-on necks.
20
+
21
+ Several other companies also began manufacturing bass guitars during the 1950s. 1956 saw the appearance at the German trade fair "Musikmesse Frankfurt" of the distinctive Höfner 500/1 violin-shaped bass, made using violin construction techniques by Walter Höfner, a second-generation violin luthier.[citation needed] The design became known as the "Beatle bass" for its use by Beatles bassist Paul McCartney. In 1957, Rickenbacker introduced the model 4000, the first bass to feature a neck-through-body design in which the neck is part of the body wood.[citation needed] Kay Musical Instrument Company began production of the K-162 in 1952, Danelectro released the Longhorn in 1956, and Burns London/Supersound in 1958.[17]
22
+
23
+ With the explosion in popularity of rock music in the 1960s, many more manufacturers began making electric basses, including Yamaha, Teisco and Guyatone. Introduced in 1960, the Fender Jazz Bass, initially known as the "Deluxe Bass", used a body design known as an offset waist which was first seen on the Jazzmaster guitar in an effort to improve comfort while playing seated.[18] The "J-bass" featured two single-coil pickups, one close to the bridge and one in the Precision bass's split coil pickup position. The earliest production Jazz basses had a pair of concentric (or "stacked") knobs to control volume and tone for each pickup; this was soon changed to the present configuration of a volume control for each pickup, and a single passive tone control.
24
+
25
+ The Jazz Bass's neck was narrower at the nut than the Precision bass — 1 1⁄2 inches (38 mm) versus 1 3⁄4 inches (44 mm) — allowing for easier access to the lower strings and an overall spacing and feel closer to that of an electric guitar, allowing trained guitarists to transition to the bass guitar more easily.[citation needed] Another visual difference that set the Jazz Bass apart from the Precision is its "offset-waist" body.[further explanation needed]
26
+
27
+ Pickup shapes on electric basses are often referred to as "P" or "J" pickups in reference to the visual and electrical differences between the Precision Bass and Jazz Bass pickups.[citation needed] In the 1950s and 1960s, all bass guitars were often called the "Fender bass", due to Fender's early dominance in the market.
28
+
29
+ Providing a more "Gibson-scale" instrument, rather than the 34 inches (864 mm) Jazz and Precision, Fender produced the Mustang Bass, a 30-inch (762 mm) scale-length instrument.[citation needed] The Fender VI, a 6 string bass, was tuned one octave lower than standard guitar tuning. It was released in 1961, and was briefly favored by Jack Bruce of Cream.[citation needed]
30
+
31
+ Gibson introduced its short-scale 30 1⁄2-inch (775 mm) EB-3 in 1961, also used by Bruce.[19] The EB-3 had a "mini-humbucker" at the bridge position. Gibson basses tended to be smaller, sleeker instruments with a shorter scale length than the Precision; Gibson did not produce a 34-inch (864 mm)-scale bass until 1963 with the release of the Thunderbird, which was also the first Gibson bass to use two humbucking pickups in a more traditional position, about halfway between the neck and bridge.[citation needed]
32
+
33
+ In 1971, Alembic established what became known as "boutique" or "high-end" electric bass guitars.[citation needed] These expensive, custom-tailored instruments, as used by Phil Lesh, Jack Casady, and Stanley Clarke, featured unique designs, premium hand-finished wood bodies, and innovative construction techniques such as multi-laminate neck-through-body construction and graphite necks. Alembic also pioneered the use of onboard electronics for pre-amplification and equalization.[citation needed] Active electronics increase the output of the instrument, and allow more options for controlling tonal flexibility, giving the player the ability to amplify as well as to attenuate certain frequency ranges while improving the overall frequency response (including more low-register and high-register sounds). 1973 saw the UK company Wal begin production of a their own range of active basses.[citation needed] In 1974 Music Man Instruments, founded by Tom Walker, Forrest White and Leo Fender, introduced the StingRay, the first widely produced bass with active (powered) electronics built into the instrument.[citation needed] Basses with active electronics can include a preamplifier and knobs for boosting and cutting the low and high frequencies.
34
+
35
+ In the mid-1970s, Alembic and other high-end manufacturers, such as Tobias, began offering five-string basses, with a very low "B" string.[citation needed] In 1975, bassist Anthony Jackson commissioned luthier Carl Thompson to build a six-string bass tuned (low to high) B0, E1, A1, D2, G2, C3, adding a low B string and a high C string.[20] These five- and six-string "extended-range basses" would become popular with session bassists, reducing the need for re-tuning to alternate detuned configurations like "drop D", and also allowing the bassist to play more notes from the same fretting position with fewer shifts up and down the fingerboard, a crucial benefit for a session player sightreading basslines at a recording session.[citation needed]
36
+
37
+ In the 1980s, bass designers continued to explore new approaches. Ned Steinberger introduced a headless bass in 1979 and continued his innovations in the 1980s, using graphite and other new materials and (in 1984) introducing the TransTrem tremolo bar. In 1982, Hans-Peter Wilfer founded Warwick, to make a European bass, as the market at the time was dominated by Asian and American basses. Their first bass was the Streamer Bass, which is similar to the Spector NS. In 1987, the Guild Guitar Corporation launched the fretless Ashbory bass, which used silicone rubber strings and a piezoelectric pickup to achieve an "upright bass" sound with a short 18-inch (457 mm) scale length. In the late 1980s, MTV's "Unplugged" show, which featured bands performing with acoustic instruments, helped to popularize hollow-bodied acoustic bass guitars amplified with piezoelectric pickups built into the bridge of the instrument.[citation needed]
38
+
39
+ During the 1990s, as five-string basses became more widely available and more affordable, an increasing number of bassists in genres ranging from metal to gospel began using five-string instruments for added lower range—a low "B" string. As well, onboard battery-powered electronics such as preamplifiers and equalizer circuits, which were previously only available on expensive "boutique" instruments, became increasingly available on mid-priced basses. From 2000 to the 2010s, some bass manufacturers included digital modelling circuits inside the instrument on more costly instruments to recreate tones and sounds from many models of basses (e.g., Line 6's Variax bass). A modelling bass can digitally emulate the tone and sound of many famous basses, ranging from a vintage Fender Precision to a Rickenbacker. However, as with the electric guitar, traditional "passive" bass designs, which include only pickups, tone and volume knobs (without a preamp or other electronics) remained popular. Reissued versions of vintage instruments such as the Fender Precision Bass and Fender Jazz Bass remained popular among new instrument buyers up to the 2010s. In 2011, a 60th Anniversary P-bass was introduced by Fender, along with the re-introduction of the short-scale Fender Jaguar Bass.[citation needed]
40
+
41
+ While electric bass guitars are traditionally fretted instruments, fretless basses are used by some players to achieve different tones. Rolling Stones bassist Bill Wyman is sometimes identified as the first to make a fretless bass. In 1961, he converted a used UK-built Dallas Tuxedo bass by removing the frets and filling in the slots with wood putty.[13] Wyman used it to record songs such as "Paint It, Black" and "Mother's Little Helper" in 1966.
42
+
43
+ In 1966, Ampeg introduced the AUB-1, the first production fretless bass. Fender followed with a fretless Precision Bass in 1970. Some fretless basses have "fret line" markers inlaid in the fingerboard as a guide, while others only use guide marks on the side of the neck. In the early 1970s, fusion-jazz bassist Jaco Pastorius coated the fingerboard of his de-fretted Fender Jazz Bass in epoxy resin, allowing him to use roundwound strings for a brighter sound without damaging the fretboard.
44
+
45
+ Traditional electric bass guitars have four strings, tuned the same as double basses: E1–A1–D2–G2. However, now there are many options, with five-, six-, and more string designs, with many approaches to tuning. In addition to traditional flatwound strings, choices now include various windings and materials.
46
+
47
+ The use of non-magnetic pickups allows bassists to use non-ferrous strings such as nylon, brass, polyurethane and silicone rubber. These materials produce different tones and, in the case of the polyurethane or silicone rubber strings, allow much shorter scale lengths.
48
+
49
+ Similar to the electric guitar, the typical electric bass guitar requires an external amplifier in order to be heard in performance settings. Additionally, various electronic effects, such as preamplifiers, "stomp box"-style pedals and signal processors are available to allow for further shaping of the sound.
en/5810.html.txt ADDED
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1
+ The tropics is the region of the Earth surrounding the Equator. They are delimited in latitude by the Tropic of Cancer in the Northern Hemisphere at 23°26′11.8″ (or 23.43662°) N and the Tropic of Capricorn in
2
+ the Southern Hemisphere at 23°26′11.8″ (or 23.43662°) S; these latitudes correspond to the axial tilt of the Earth. The tropics are also referred to as the tropical zone and the torrid zone (see geographical zone). The tropics include all the areas on the Earth where the Sun contacts a point directly overhead at least once during the solar year (which is a subsolar point) - thus the latitude of the tropics is roughly equal to the angle of the Earth's axial tilt.
3
+
4
+ In terms of climate, the tropics receive sunlight that is more direct than the rest of Earth and are generally hotter and wetter. The word "tropical" sometimes refers to this sort of climate rather than to the geographical zone. The tropical zone includes deserts and snow-capped mountains, which are not tropical in the climatic sense. The tropics are distinguished from the other climatic and biomatic regions of Earth, which are the middle latitudes and the polar regions on either side of the equatorial zone.
5
+
6
+ The tropics constitute 40% of the Earth's surface area[1] and contain 36% of the Earth's landmass.[2] As of 2014[update], the region was home to 40% of the world's population, and this figure was then projected to reach 50% by 2050.[3]
7
+
8
+ The word "tropic" comes from Ancient Greek τροπή (tropē), meaning "to turn" or "change direction", as the Sun appears to cease its southerly course when it reaches this latitude and begins moving back to the north.
9
+
10
+ "Tropical" is sometimes used in a general sense for a tropical climate to mean warm to hot and moist year-round, often with the sense of lush vegetation.
11
+
12
+ Many tropical areas have a dry and wet season. The wet season, rainy season or green season is the time of year, ranging from one or more months, when most of the average annual rainfall in a region falls.[4] Areas with wet seasons are disseminated across portions of the tropics and subtropics.[5] Under the Köppen climate classification, for tropical climates, a wet-season month is defined as a month where average precipitation is 60 millimetres (2.4 in) or more.[6] Tropical rainforests technically do not have dry or wet seasons, since their rainfall is equally distributed through the year.[7] Some areas with pronounced rainy seasons see a break in rainfall during mid-season when the intertropical convergence zone or monsoon trough moves poleward of their location during the middle of the warm season;[8] typical vegetation in these areas ranges from moist seasonal tropical forests to savannahs.
13
+
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+ When the wet season occurs during the warm season, or summer, precipitation falls mainly during the late afternoon and early evening hours. The wet season is a time when air quality improves, freshwater quality improves and vegetation grows significantly, leading to crop yields late in the season. Floods cause rivers to overflow their banks, and some animals to retreat to higher ground. Soil nutrients diminish and erosion increases. The incidence of malaria increases in areas where the rainy season coincides with high temperatures. Animals have adaptation and survival strategies for the wetter regime. The previous dry season leads to food shortages into the wet season, as the crops have yet to mature.
15
+
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+ However, regions within the tropics may well not have a tropical climate. Under the Köppen climate classification, much of the area within the geographical tropics is classed not as "tropical" but as "dry" (arid or semi-arid), including the Sahara Desert, the Atacama Desert and Australian Outback. Also, there are alpine tundra and snow-capped peaks, including Mauna Kea, Mount Kilimanjaro, and the Andes as far south as the northernmost parts of Chile and Perú.
17
+
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+ Tropical plants and animals are those species native to the tropics. Tropical ecosystems may consist of tropical rainforests, seasonal tropical forests, dry (often deciduous) forests, spiny forests, desert and other habitat types. There are often significant areas of biodiversity, and species endemism present, particularly in rainforests and seasonal forests. Some examples of important biodiversity and high endemism ecosystems are El Yunque National Forest in Puerto Rico, Costa Rican and Nicaraguan rainforests, Amazon Rainforest territories of several South American countries, Madagascar dry deciduous forests, the Waterberg Biosphere of South Africa, and eastern Madagascar rainforests. Often the soils of tropical forests are low in nutrient content, making them quite vulnerable to slash-and-burn deforestation techniques, which are sometimes an element of shifting cultivation agricultural systems.
19
+
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+ In biogeography, the tropics are divided into Paleotropics (Africa, Asia and Australia) and Neotropics (Caribbean, Central America, and South America). Together, they are sometimes referred to as the Pantropic. The system of biogeographic realms differs somewhat; the Neotropical realm includes both the Neotropics and temperate South America, and the Paleotropics correspond to the Afrotropical, Indomalayan, Oceanian, and tropical Australasian realms.
21
+
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+ Tropicality refers to the image that people outside the tropics have of the region, ranging from critical to verging on fetishism. The idea of tropicality gained renewed interest in geographical discourse when French geographer Pierre Gourou published Les Pays Tropicaux (The Tropical World in English), in the late 1940s.[9]
23
+
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+ Tropicality encompassed two images. One, is that the tropics represent a 'Garden of Eden', a heaven on Earth, a land of rich biodiversity - aka a tropical paradise.[10] The alternative is that the tropics consist of wild, unconquerable nature. The latter view was often discussed in old Western literature more so than the first.[10] Evidence suggests over time that the view of the tropics as such in popular literature has been supplanted by more well-rounded and sophisticated interpretations.[11]
25
+
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+ Western scholars tried to theorize reasons about why tropical areas were relatively more inhospitable to human civilisations then those existing in colder regions of the Northern Hemisphere. A popular explanation focused on the differences in climate. Tropical jungles and rainforests have much more humid and hotter weather than colder and drier temperaments of the Northern Hemisphere. This theme led to some scholars to suggest that humid hot climates correlate to human populations lacking control over nature e.g. ' the wild Amazonian rainforests'.[12]
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1
+
2
+
3
+ A black hole is a region of spacetime where gravity is so strong that nothing—no particles or even electromagnetic radiation such as light—can escape from it.[6] The theory of general relativity predicts that a sufficiently compact mass can deform spacetime to form a black hole.[7][8] The boundary of the region from which no escape is possible is called the event horizon. Although the event horizon has an enormous effect on the fate and circumstances of an object crossing it, according to general relativity it has no locally detectable features.[9] In many ways, a black hole acts like an ideal black body, as it reflects no light.[10][11] Moreover, quantum field theory in curved spacetime predicts that event horizons emit Hawking radiation, with the same spectrum as a black body of a temperature inversely proportional to its mass. This temperature is on the order of billionths of a kelvin for black holes of stellar mass, making it essentially impossible to observe.
4
+
5
+ Objects whose gravitational fields are too strong for light to escape were first considered in the 18th century by John Michell and Pierre-Simon Laplace.[12] The first modern solution of general relativity that would characterize a black hole was found by Karl Schwarzschild in 1916, although its interpretation as a region of space from which nothing can escape was first published by David Finkelstein in 1958. Black holes were long considered a mathematical curiosity; it was not until the 1960s that theoretical work showed they were a generic prediction of general relativity. The discovery of neutron stars by Jocelyn Bell Burnell in 1967 sparked interest in gravitationally collapsed compact objects as a possible astrophysical reality.
6
+
7
+ Black holes of stellar mass are expected to form when very massive stars collapse at the end of their life cycle. After a black hole has formed, it can continue to grow by absorbing mass from its surroundings. By absorbing other stars and merging with other black holes, supermassive black holes of millions of solar masses (M☉) may form. There is consensus that supermassive black holes exist in the centers of most galaxies.
8
+
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+ The presence of a black hole can be inferred through its interaction with other matter and with electromagnetic radiation such as visible light. Matter that falls onto a black hole can form an external accretion disk heated by friction, forming quasars, some of the brightest objects in the universe. Stars passing too close to a supermassive black hole can be shred into streamers that shine very brightly before being "swallowed."[13] If there are other stars orbiting a black hole, their orbits can be used to determine the black hole's mass and location. Such observations can be used to exclude possible alternatives such as neutron stars. In this way, astronomers have identified numerous stellar black hole candidates in binary systems, and established that the radio source known as Sagittarius A*, at the core of the Milky Way galaxy, contains a supermassive black hole of about 4.3 million solar masses.
10
+
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+ On 11 February 2016, the LIGO Scientific Collaboration and the Virgo collaboration announced the first direct detection of gravitational waves, which also represented the first observation of a black hole merger.[14] As of December 2018[update], eleven gravitational wave events have been observed that originated from ten merging black holes (along with one binary neutron star merger).[15][16] On 10 April 2019, the first direct image of a black hole and its vicinity was published, following observations made by the Event Horizon Telescope in 2017 of the supermassive black hole in Messier 87's galactic centre.[3][17][18]
12
+
13
+ The idea of a body so massive that even light could not escape was briefly proposed by astronomical pioneer and English clergyman John Michell in a letter published in November 1784. Michell's simplistic calculations assumed such a body might have the same density as the Sun, and concluded that such a body would form when a star's diameter exceeds the Sun's by a factor of 500, and the surface escape velocity exceeds the usual speed of light. Michell correctly noted that such supermassive but non-radiating bodies might be detectable through their gravitational effects on nearby visible bodies.[20][12][21] Scholars of the time were initially excited by the proposal that giant but invisible stars might be hiding in plain view, but enthusiasm dampened when the wavelike nature of light became apparent in the early nineteenth century.[22]
14
+
15
+ If light were a wave rather than a "corpuscle", it is unclear what, if any, influence gravity would have on escaping light waves.[12][21] Modern physics discredits Michell's notion of a light ray shooting directly from the surface of a supermassive star, being slowed down by the star's gravity, stopping, and then free-falling back to the star's surface.[23]
16
+
17
+ In 1915, Albert Einstein developed his theory of general relativity, having earlier shown that gravity does influence light's motion. Only a few months later, Karl Schwarzschild found a solution to the Einstein field equations, which describes the gravitational field of a point mass and a spherical mass.[24] A few months after Schwarzschild, Johannes Droste, a student of Hendrik Lorentz, independently gave the same solution for the point mass and wrote more extensively about its properties.[25][26] This solution had a peculiar behaviour at what is now called the Schwarzschild radius, where it became singular, meaning that some of the terms in the Einstein equations became infinite. The nature of this surface was not quite understood at the time. In 1924, Arthur Eddington showed that the singularity disappeared after a change of coordinates (see Eddington–Finkelstein coordinates), although it took until 1933 for Georges Lemaître to realize that this meant the singularity at the Schwarzschild radius was a non-physical coordinate singularity.[27] Arthur Eddington did however comment on the possibility of a star with mass compressed to the Schwarzschild radius in a 1926 book, noting that Einstein's theory allows us to rule out overly large densities for visible stars like Betelgeuse because "a star of 250 million km radius could not possibly have so high a density as the sun. Firstly, the force of gravitation would be so great that light would be unable to escape from it, the rays falling back to the star like a stone to the earth. Secondly, the red shift of the spectral lines would be so great that the spectrum would be shifted out of existence. Thirdly, the mass would produce so much curvature of the space-time metric that space would close up around the star, leaving us outside (i.e., nowhere)."[28][29]
18
+
19
+ In 1931, Subrahmanyan Chandrasekhar calculated, using special relativity, that a non-rotating body of electron-degenerate matter above a certain limiting mass (now called the Chandrasekhar limit at 1.4 M☉) has no stable solutions.[30] His arguments were opposed by many of his contemporaries like Eddington and Lev Landau, who argued that some yet unknown mechanism would stop the collapse.[31] They were partly correct: a white dwarf slightly more massive than the Chandrasekhar limit will collapse into a neutron star,[32] which is itself stable. But in 1939, Robert Oppenheimer and others predicted that neutron stars above another limit (the Tolman–Oppenheimer–Volkoff limit) would collapse further for the reasons presented by Chandrasekhar, and concluded that no law of physics was likely to intervene and stop at least some stars from collapsing to black holes.[33] Their original calculations, based on the Pauli exclusion principle, gave it as 0.7 M☉; subsequent consideration of strong force-mediated neutron-neutron repulsion raised the estimate to approximately 1.5 M☉ to 3.0 M☉.[34] Observations of the neutron star merger GW170817, which is thought to have generated a black hole shortly afterward, have refined the TOV limit estimate to ~2.17 M☉.[35][36][37][38][39]
20
+
21
+ Oppenheimer and his co-authors interpreted the singularity at the boundary of the Schwarzschild radius as indicating that this was the boundary of a bubble in which time stopped. This is a valid point of view for external observers, but not for infalling observers. Because of this property, the collapsed stars were called "frozen stars", because an outside observer would see the surface of the star frozen in time at the instant where its collapse takes it to the Schwarzschild radius.[40]
22
+
23
+ In 1958, David Finkelstein identified the Schwarzschild surface as an event horizon, "a perfect unidirectional membrane: causal influences can cross it in only one direction".[41] This did not strictly contradict Oppenheimer's results, but extended them to include the point of view of infalling observers. Finkelstein's solution extended the Schwarzschild solution for the future of observers falling into a black hole. A complete extension had already been found by Martin Kruskal, who was urged to publish it.[42]
24
+
25
+ These results came at the beginning of the golden age of general relativity, which was marked by general relativity and black holes becoming mainstream subjects of research. This process was helped by the discovery of pulsars by Jocelyn Bell Burnell in 1967,[43][44] which, by 1969, were shown to be rapidly rotating neutron stars.[45] Until that time, neutron stars, like black holes, were regarded as just theoretical curiosities; but the discovery of pulsars showed their physical relevance and spurred a further interest in all types of compact objects that might be formed by gravitational collapse.[citation needed]
26
+
27
+ In this period more general black hole solutions were found. In 1963, Roy Kerr found the exact solution for a rotating black hole. Two years later, Ezra Newman found the axisymmetric solution for a black hole that is both rotating and electrically charged.[46] Through the work of Werner Israel,[47] Brandon Carter,[48][49] and David Robinson[50] the no-hair theorem emerged, stating that a stationary black hole solution is completely described by the three parameters of the Kerr–Newman metric: mass, angular momentum, and electric charge.[51]
28
+
29
+ At first, it was suspected that the strange features of the black hole solutions were pathological artifacts from the symmetry conditions imposed, and that the singularities would not appear in generic situations. This view was held in particular by Vladimir Belinsky, Isaak Khalatnikov, and Evgeny Lifshitz, who tried to prove that no singularities appear in generic solutions. However, in the late 1960s Roger Penrose[52] and Stephen Hawking used global techniques to prove that singularities appear generically.[53]
30
+
31
+ Work by James Bardeen, Jacob Bekenstein, Carter, and Hawking in the early 1970s led to the formulation of black hole thermodynamics.[54] These laws describe the behaviour of a black hole in close analogy to the laws of thermodynamics by relating mass to energy, area to entropy, and surface gravity to temperature. The analogy was completed when Hawking, in 1974, showed that quantum field theory implies that black holes should radiate like a black body with a temperature proportional to the surface gravity of the black hole, predicting the effect now known as Hawking radiation.[55]
32
+
33
+ John Michell used the term "dark star",[56] and in the early 20th century, physicists used the term "gravitationally collapsed object". Science writer Marcia Bartusiak traces the term "black hole" to physicist Robert H. Dicke, who in the early 1960s reportedly compared the phenomenon to the Black Hole of Calcutta, notorious as a prison where people entered but never left alive.[57]
34
+
35
+ The term "black hole" was used in print by Life and Science News magazines in 1963,[57] and by science journalist Ann Ewing in her article "'Black Holes' in Space", dated 18 January 1964, which was a report on a meeting of the American Association for the Advancement of Science held in Cleveland, Ohio.[58][59]
36
+
37
+ In December 1967, a student reportedly suggested the phrase "black hole" at a lecture by John Wheeler;[58] Wheeler adopted the term for its brevity and "advertising value", and it quickly caught on,[60] leading some to credit Wheeler with coining the phrase.[61]
38
+
39
+ The no-hair conjecture postulates that, once it achieves a stable condition after formation, a black hole has only three independent physical properties: mass, charge, and angular momentum; the black hole is otherwise featureless. If the conjecture is true, any two black holes that share the same values for these properties, or parameters, are indistinguishable from one another. The degree to which the conjecture is true for real black holes under the laws of modern physics, is currently an unsolved problem.[51]
40
+
41
+ These properties are special because they are visible from outside a black hole. For example, a charged black hole repels other like charges just like any other charged object. Similarly, the total mass inside a sphere containing a black hole can be found by using the gravitational analog of Gauss's law, the ADM mass, far away from the black hole.[62][clarification needed] Likewise, the angular momentum can be measured from far away using frame dragging by the gravitomagnetic field.[clarification needed]
42
+
43
+ When an object falls into a black hole, any information about the shape of the object or distribution of charge on it is evenly distributed along the horizon of the black hole, and is lost to outside observers. The behavior of the horizon in this situation is a dissipative system that is closely analogous to that of a conductive stretchy membrane with friction and electrical resistance—the membrane paradigm.[63] This is different from other field theories such as electromagnetism, which do not have any friction or resistivity at the microscopic level, because they are time-reversible. Because a black hole eventually achieves a stable state with only three parameters, there is no way to avoid losing information about the initial conditions: the gravitational and electric fields of a black hole give very little information about what went in. The information that is lost includes every quantity that cannot be measured far away from the black hole horizon, including approximately conserved quantum numbers such as the total baryon number and lepton number. This behavior is so puzzling that it has been called the black hole information loss paradox.[64][65]
44
+
45
+ The simplest static black holes have mass but neither electric charge nor angular momentum. These black holes are often referred to as Schwarzschild black holes after Karl Schwarzschild who discovered this solution in 1916.[24] According to Birkhoff's theorem, it is the only vacuum solution that is spherically symmetric.[66] This means there is no observable difference at a distance between the gravitational field of such a black hole and that of any other spherical object of the same mass. The popular notion of a black hole "sucking in everything" in its surroundings is therefore correct only near a black hole's horizon; far away, the external gravitational field is identical to that of any other body of the same mass.[67]
46
+
47
+ Solutions describing more general black holes also exist. Non-rotating charged black holes are described by the Reissner–Nordström metric, while the Kerr metric describes a non-charged rotating black hole. The most general stationary black hole solution known is the Kerr–Newman metric, which describes a black hole with both charge and angular momentum.[68]
48
+
49
+ While the mass of a black hole can take any positive value, the charge and angular momentum are constrained by the mass. In Planck units, the total electric charge Q and the total angular momentum J are expected to satisfy
50
+
51
+ for a black hole of mass M. Black holes with the minimum possible mass satisfying this inequality are called extremal. Solutions of Einstein's equations that violate this inequality exist, but they do not possess an event horizon. These solutions have so-called naked singularities that can be observed from the outside, and hence are deemed unphysical. The cosmic censorship hypothesis rules out the formation of such singularities, when they are created through the gravitational collapse of realistic matter.[7] This is supported by numerical simulations.[69]
52
+
53
+ Due to the relatively large strength of the electromagnetic force, black holes forming from the collapse of stars are expected to retain the nearly neutral charge of the star. Rotation, however, is expected to be a universal feature of compact astrophysical objects. The black-hole candidate binary X-ray source GRS 1915+105[70] appears to have an angular momentum near the maximum allowed value. That uncharged limit is[71]
54
+
55
+ allowing definition of a dimensionless spin parameter such that[71]
56
+
57
+ Black holes are commonly classified according to their mass, independent of angular momentum, J. The size of a black hole, as determined by the radius of the event horizon, or Schwarzschild radius, is proportional to the mass, M, through
58
+
59
+ where rs is the Schwarzschild radius and MSun is the mass of the Sun.[73] For a black hole with nonzero spin and/or electric charge, the radius is smaller,[Note 2] until an extremal black hole could have an event horizon close to[74]
60
+
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+ The defining feature of a black hole is the appearance of an event horizon—a boundary in spacetime through which matter and light can pass only inward towards the mass of the black hole. Nothing, not even light, can escape from inside the event horizon. The event horizon is referred to as such because if an event occurs within the boundary, information from that event cannot reach an outside observer, making it impossible to determine whether such an event occurred.[76]
62
+
63
+ As predicted by general relativity, the presence of a mass deforms spacetime in such a way that the paths taken by particles bend towards the mass.[77] At the event horizon of a black hole, this deformation becomes so strong that there are no paths that lead away from the black hole.[78]
64
+
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+ To a distant observer, clocks near a black hole would appear to tick more slowly than those further away from the black hole.[79] Due to this effect, known as gravitational time dilation, an object falling into a black hole appears to slow as it approaches the event horizon, taking an infinite time to reach it.[80] At the same time, all processes on this object slow down, from the view point of a fixed outside observer, causing any light emitted by the object to appear redder and dimmer, an effect known as gravitational redshift.[81] Eventually, the falling object fades away until it can no longer be seen. Typically this process happens very rapidly with an object disappearing from view within less than a second.[82]
66
+
67
+ On the other hand, indestructible observers falling into a black hole do not notice any of these effects as they cross the event horizon. According to their own clocks, which appear to them to tick normally, they cross the event horizon after a finite time without noting any singular behaviour; in classical general relativity, it is impossible to determine the location of the event horizon from local observations, due to Einstein's equivalence principle.[83][84]
68
+
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+ The topology of the event horizon of a black hole at equilibrium is always spherical.[Note 4][87] For non-rotating (static) black holes the geometry of the event horizon is precisely spherical, while for rotating black holes the event horizon is oblate.[88][89][90]
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+
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+ At the center of a black hole, as described by general relativity, may lie a gravitational singularity, a region where the spacetime curvature becomes infinite.[91] For a non-rotating black hole, this region takes the shape of a single point and for a rotating black hole, it is smeared out to form a ring singularity that lies in the plane of rotation.[92] In both cases, the singular region has zero volume. It can also be shown that the singular region contains all the mass of the black hole solution.[93] The singular region can thus be thought of as having infinite density.[94]
72
+
73
+ Observers falling into a Schwarzschild black hole (i.e., non-rotating and not charged) cannot avoid being carried into the singularity once they cross the event horizon. They can prolong the experience by accelerating away to slow their descent, but only up to a limit.[95] When they reach the singularity, they are crushed to infinite density and their mass is added to the total of the black hole. Before that happens, they will have been torn apart by the growing tidal forces in a process sometimes referred to as spaghettification or the "noodle effect".[96]
74
+
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+ In the case of a charged (Reissner–Nordström) or rotating (Kerr) black hole, it is possible to avoid the singularity. Extending these solutions as far as possible reveals the hypothetical possibility of exiting the black hole into a different spacetime with the black hole acting as a wormhole.[97] The possibility of traveling to another universe is, however, only theoretical since any perturbation would destroy this possibility.[98] It also appears to be possible to follow closed timelike curves (returning to one's own past) around the Kerr singularity, which leads to problems with causality like the grandfather paradox.[99] It is expected that none of these peculiar effects would survive in a proper quantum treatment of rotating and charged black holes.[100]
76
+
77
+ The appearance of singularities in general relativity is commonly perceived as signaling the breakdown of the theory.[101] This breakdown, however, is expected; it occurs in a situation where quantum effects should describe these actions, due to the extremely high density and therefore particle interactions. To date, it has not been possible to combine quantum and gravitational effects into a single theory, although there exist attempts to formulate such a theory of quantum gravity. It is generally expected that such a theory will not feature any singularities.[102][103]
78
+
79
+ The photon sphere is a spherical boundary of zero thickness in which photons that move on tangents to that sphere would be trapped in a circular orbit about the black hole. For non-rotating black holes, the photon sphere has a radius 1.5 times the Schwarzschild radius. Their orbits would be dynamically unstable, hence any small perturbation, such as a particle of infalling matter, would cause an instability that would grow over time, either setting the photon on an outward trajectory causing it to escape the black hole, or on an inward spiral where it would eventually cross the event horizon.[104]
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+
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+ While light can still escape from the photon sphere, any light that crosses the photon sphere on an inbound trajectory will be captured by the black hole. Hence any light that reaches an outside observer from the photon sphere must have been emitted by objects between the photon sphere and the event horizon.[104]
82
+
83
+ Rotating black holes are surrounded by a region of spacetime in which it is impossible to stand still, called the ergosphere. This is the result of a process known as frame-dragging; general relativity predicts that any rotating mass will tend to slightly "drag" along the spacetime immediately surrounding it. Any object near the rotating mass will tend to start moving in the direction of rotation. For a rotating black hole, this effect is so strong near the event horizon that an object would have to move faster than the speed of light in the opposite direction to just stand still.[106]
84
+
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+ The ergosphere of a black hole is a volume whose inner boundary is the black hole's event horizon and an outer boundary called the ergosurface, which coincides with the event horizon at the poles but noticeably wider around the equator.[105]
86
+
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+ Objects and radiation can escape normally from the ergosphere. Through the Penrose process, objects can emerge from the ergosphere with more energy than they entered. This energy is taken from the rotational energy of the black hole causing the latter to slow.[107] A variation of the Penrose process in the presence of strong magnetic fields, the Blandford–Znajek process is considered a likely mechanism for the enormous luminosity and relativistic jets of quasars and other active galactic nuclei.
88
+
89
+ In Newtonian gravity, test particles can stably orbit at arbitrary distances from a central object. In general relativity, however, there exists an innermost stable circular orbit (often called the ISCO), inside of which, any infinitesimal perturbations to a circular orbit will lead to inspiral into the black hole.[108] The location of the ISCO depends on the spin of the black hole, in the case of a Schwarzschild black hole (spin zero) is:
90
+
91
+ and decreases with increasing black hole spin for particles orbiting in the same direction as the spin.[109]
92
+
93
+ Given the bizarre character of black holes, it was long questioned whether such objects could actually exist in nature or whether they were merely pathological solutions to Einstein's equations. Einstein himself wrongly thought black holes would not form, because he held that the angular momentum of collapsing particles would stabilize their motion at some radius.[110] This led the general relativity community to dismiss all results to the contrary for many years. However, a minority of relativists continued to contend that black holes were physical objects,[111] and by the end of the 1960s, they had persuaded the majority of researchers in the field that there is no obstacle to the formation of an event horizon.[citation needed]
94
+
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+ Penrose demonstrated that once an event horizon forms, general relativity without quantum mechanics requires that a singularity will form within.[52] Shortly afterwards, Hawking showed that many cosmological solutions that describe the Big Bang have singularities without scalar fields or other exotic matter (see "Penrose–Hawking singularity theorems").[clarification needed] The Kerr solution, the no-hair theorem, and the laws of black hole thermodynamics showed that the physical properties of black holes were simple and comprehensible, making them respectable subjects for research.[112] Conventional black holes are formed by gravitational collapse of heavy objects such as stars, but they can also in theory be formed by other processes.[113][114]
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+
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+ Gravitational collapse occurs when an object's internal pressure is insufficient to resist the object's own gravity. For stars this usually occurs either because a star has too little "fuel" left to maintain its temperature through stellar nucleosynthesis, or because a star that would have been stable receives extra matter in a way that does not raise its core temperature. In either case the star's temperature is no longer high enough to prevent it from collapsing under its own weight.[115]
98
+ The collapse may be stopped by the degeneracy pressure of the star's constituents, allowing the condensation of matter into an exotic denser state. The result is one of the various types of compact star. Which type forms depends on the mass of the remnant of the original star left after the outer layers have been blown away. Such explosions and pulsations lead to planetary nebula.[116] This mass can be substantially less than the original star. Remnants exceeding 5 M☉ are produced by stars that were over 20 M☉ before the collapse.[115]
99
+
100
+ If the mass of the remnant exceeds about 3–4 M☉ (the Tolman–Oppenheimer–Volkoff limit[33]), either because the original star was very heavy or because the remnant collected additional mass through accretion of matter, even the degeneracy pressure of neutrons is insufficient to stop the collapse. No known mechanism (except possibly quark degeneracy pressure, see quark star) is powerful enough to stop the implosion and the object will inevitably collapse to form a black hole.[115]
101
+
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+ The gravitational collapse of heavy stars is assumed to be responsible for the formation of stellar mass black holes. Star formation in the early universe may have resulted in very massive stars, which upon their collapse would have produced black holes of up to 103 M☉. These black holes could be the seeds of the supermassive black holes found in the centers of most galaxies.[118] It has further been suggested that supermassive black holes with typical masses of ~105 M☉ could have formed from the direct collapse of gas clouds in the young universe.[113] Some candidates for such objects have been found in observations of the young universe.[113]
103
+
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+ While most of the energy released during gravitational collapse is emitted very quickly, an outside observer does not actually see the end of this process. Even though the collapse takes a finite amount of time from the reference frame of infalling matter, a distant observer would see the infalling material slow and halt just above the event horizon, due to gravitational time dilation. Light from the collapsing material takes longer and longer to reach the observer, with the light emitted just before the event horizon forms delayed an infinite amount of time. Thus the external observer never sees the formation of the event horizon; instead, the collapsing material seems to become dimmer and increasingly red-shifted, eventually fading away.[119]
105
+
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+ Gravitational collapse requires great density. In the current epoch of the universe these high densities are found only in stars, but in the early universe shortly after the Big Bang densities were much greater, possibly allowing for the creation of black holes. High density alone is not enough to allow black hole formation since a uniform mass distribution will not allow the mass to bunch up. In order for primordial black holes to have formed in such a dense medium, there must have been initial density perturbations that could then grow under their own gravity. Different models for the early universe vary widely in their predictions of the scale of these fluctuations. Various models predict the creation of primordial black holes ranging in size from a Planck mass to hundreds of thousands of solar masses.[114]
107
+
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+ Despite the early universe being extremely dense—far denser than is usually required to form a black hole—it did not re-collapse into a black hole during the Big Bang. Models for gravitational collapse of objects of relatively constant size, such as stars, do not necessarily apply in the same way to rapidly expanding space such as the Big Bang.[120]
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+
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+ Gravitational collapse is not the only process that could create black holes. In principle, black holes could be formed in high-energy collisions that achieve sufficient density. As of 2002, no such events have been detected, either directly or indirectly as a deficiency of the mass balance in particle accelerator experiments.[121] This suggests that there must be a lower limit for the mass of black holes. Theoretically, this boundary is expected to lie around the Planck mass (mP=√ħ c/G ≈ 1.2×1019 GeV/c2 ≈ 2.2×10−8 kg), where quantum effects are expected to invalidate the predictions of general relativity.[122] This would put the creation of black holes firmly out of reach of any high-energy process occurring on or near the Earth. However, certain developments in quantum gravity suggest that the minimum black hole mass could be much lower: some braneworld scenarios for example put the boundary as low as 1 TeV/c2.[123] This would make it conceivable for micro black holes to be created in the high-energy collisions that occur when cosmic rays hit the Earth's atmosphere, or possibly in the Large Hadron Collider at CERN. These theories are very speculative, and the creation of black holes in these processes is deemed unlikely by many specialists.[124] Even if micro black holes could be formed, it is expected that they would evaporate in about 10−25 seconds, posing no threat to the Earth.[125]
111
+
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+ Once a black hole has formed, it can continue to grow by absorbing additional matter. Any black hole will continually absorb gas and interstellar dust from its surroundings. This is the primary process through which supermassive black holes seem to have grown.[118] A similar process has been suggested for the formation of intermediate-mass black holes found in globular clusters.[126] Black holes can also merge with other objects such as stars or even other black holes. This is thought to have been important, especially in the early growth of supermassive black holes, which could have formed from the aggregation of many smaller objects.[118] The process has also been proposed as the origin of some intermediate-mass black holes.[127][128]
113
+
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+ In 1974, Hawking predicted that black holes are not entirely black but emit small amounts of thermal radiation at a temperature ℏ c3/(8 π G M kB);[55] this effect has become known as Hawking radiation. By applying quantum field theory to a static black hole background, he determined that a black hole should emit particles that display a perfect black body spectrum. Since Hawking's publication, many others have verified the result through various approaches.[129] If Hawking's theory of black hole radiation is correct, then black holes are expected to shrink and evaporate over time as they lose mass by the emission of photons and other particles.[55] The temperature of this thermal spectrum (Hawking temperature) is proportional to the surface gravity of the black hole, which, for a Schwarzschild black hole, is inversely proportional to the mass. Hence, large black holes emit less radiation than small black holes.[130]
115
+
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+ A stellar black hole of 1 M☉ has a Hawking temperature of 62 nanokelvins.[131] This is far less than the 2.7 K temperature of the cosmic microwave background radiation. Stellar-mass or larger black holes receive more mass from the cosmic microwave background than they emit through Hawking radiation and thus will grow instead of shrinking.[132] To have a Hawking temperature larger than 2.7 K (and be able to evaporate), a black hole would need a mass less than the Moon. Such a black hole would have a diameter of less than a tenth of a millimeter.[133]
117
+
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+ If a black hole is very small, the radiation effects are expected to become very strong. A black hole with the mass of a car would have a diameter of about 10−24 m and take a nanosecond to evaporate, during which time it would briefly have a luminosity of more than 200 times that of the Sun. Lower-mass black holes are expected to evaporate even faster; for example, a black hole of mass 1 TeV/c2 would take less than 10−88 seconds to evaporate completely. For such a small black hole, quantum gravitation effects are expected to play an important role and could hypothetically make such a small black hole stable, although current developments in quantum gravity do not indicate this is the case.[134][135]
119
+
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+ The Hawking radiation for an astrophysical black hole is predicted to be very weak and would thus be exceedingly difficult to detect from Earth. A possible exception, however, is the burst of gamma rays emitted in the last stage of the evaporation of primordial black holes. Searches for such flashes have proven unsuccessful and provide stringent limits on the possibility of existence of low mass primordial black holes.[136] NASA's Fermi Gamma-ray Space Telescope launched in 2008 will continue the search for these flashes.[137]
121
+
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+ If black holes evaporate via Hawking radiation, a solar mass black hole will evaporate (beginning once the temperature of the cosmic microwave background drops below that of the black hole) over a period of 1064 years.[138] A supermassive black hole with a mass of 1011 (100 billion) M☉ will evaporate in around 2×10100 years.[139] Some monster black holes in the universe are predicted to continue to grow up to perhaps 1014 M☉ during the collapse of superclusters of galaxies. Even these would evaporate over a timescale of up to 10106 years.[138]
123
+
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+ By nature, black holes do not themselves emit any electromagnetic radiation other than the hypothetical Hawking radiation, so astrophysicists searching for black holes must generally rely on indirect observations. For example, a black hole's existence can sometimes be inferred by observing its gravitational influence upon its surroundings.[140]
125
+
126
+ On 10 April 2019 an image was released of a black hole, which is seen in magnified fashion because the light paths near the event horizon are highly bent. The dark shadow in the middle results from light paths absorbed by the black hole. The image is in false color, as the detected light halo in this image is not in the visible spectrum, but radio waves.
127
+
128
+ The Event Horizon Telescope (EHT), is an active program that directly observes the immediate environment of the event horizon of black holes, such as the black hole at the centre of the Milky Way. In April 2017, EHT began observation of the black hole in the center of Messier 87.[141] "In all, eight radio observatories on six mountains and four continents observed the galaxy in Virgo on and off for 10 days in April 2017" to provide the data yielding the image two years later in April 2019.[142] After two years of data processing, EHT released the first direct image of a black hole, specifically the supermassive black hole that lies in the center of the aforementioned galaxy.[143][144] What is visible is not the black hole, which shows as black because of the loss of all light within this dark region, rather it is the gases at the edge of the event horizon, which are displayed as orange or red, that define the black hole.[145]
129
+
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+ The brightening of this material in the 'bottom' half of the processed EHT image is thought to be caused by Doppler beaming, whereby material approaching the viewer at relativistic speeds is perceived as brighter than material moving away. In the case of a black hole this phenomenon implies that the visible material is rotating at relativistic speeds (>1,000 km/s), the only speeds at which it is possible to centrifugally balance the immense gravitational attraction of the singularity, and thereby remain in orbit above the event horizon. This configuration of bright material implies that the EHT observed M87* from a perspective catching the black hole's accretion disc nearly edge-on, as the whole system rotated clockwise.[146] However, the extreme gravitational lensing associated with black holes produces the illusion of a perspective that sees the accretion disc from above. In reality, most of the ring in the EHT image was created when the light emitted by the far side of the accretion disc bent around the black hole's gravity well and escaped such that most of the possible perspectives on M87* can see the entire disc, even that directly behind the "shadow".
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+
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+ Prior to this, in 2015, the EHT detected magnetic fields just outside the event horizon of Sagittarius A*, and even discerned some of their properties. The field lines that pass through the accretion disc were found to be a complex mixture of ordered and tangled. The existence of magnetic fields had been predicted by theoretical studies of black holes.[147][148]
133
+
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+ On 14 September 2015 the LIGO gravitational wave observatory made the first-ever successful direct observation of gravitational waves.[14][150] The signal was consistent with theoretical predictions for the gravitational waves produced by the merger of two black holes: one with about 36 solar masses, and the other around 29 solar masses.[14][151] This observation provides the most concrete evidence for the existence of black holes to date. For instance, the gravitational wave signal suggests that the separation of the two objects prior to the merger was just 350 km (or roughly four times the Schwarzschild radius corresponding to the inferred masses). The objects must therefore have been extremely compact, leaving black holes as the most plausible interpretation.[14]
135
+
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+ More importantly, the signal observed by LIGO also included the start of the post-merger ringdown, the signal produced as the newly formed compact object settles down to a stationary state. Arguably, the ringdown is the most direct way of observing a black hole.[152] From the LIGO signal it is possible to extract the frequency and damping time of the dominant mode of the ringdown. From these it is possible to infer the mass and angular momentum of the final object, which match independent predictions from numerical simulations of the merger.[153] The frequency and decay time of the dominant mode are determined by the geometry of the photon sphere. Hence, observation of this mode confirms the presence of a photon sphere, however it cannot exclude possible exotic alternatives to black holes that are compact enough to have a photon sphere.[152]
137
+
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+ The observation also provides the first observational evidence for the existence of stellar-mass black hole binaries. Furthermore, it is the first observational evidence of stellar-mass black holes weighing 25 solar masses or more.[154]
139
+
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+ Since then many more gravitational wave events have since been observed.[16]
141
+
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+ The proper motions of stars near the center of our own Milky Way provide strong observational evidence that these stars are orbiting a supermassive black hole.[155] Since 1995, astronomers have tracked the motions of 90 stars orbiting an invisible object coincident with the radio source Sagittarius A*. By fitting their motions to Keplerian orbits, the astronomers were able to infer, in 1998, that a 2.6 million M☉ object must be contained in a volume with a radius of 0.02 light-years to cause the motions of those stars.[156] Since then, one of the stars—called S2—has completed a full orbit. From the orbital data, astronomers were able to refine the calculations of the mass to 4.3 million M☉ and a radius of less than 0.002 light years for the object causing the orbital motion of those stars.[155] The upper limit on the object's size is still too large to test whether it is smaller than its Schwarzschild radius; nevertheless, these observations strongly suggest that the central object is a supermassive black hole as there are no other plausible scenarios for confining so much invisible mass into such a small volume.[156] Additionally, there is some observational evidence that this object might possess an event horizon, a feature unique to black holes.[157]
143
+
144
+ Due to conservation of angular momentum,[159] gas falling into the gravitational well created by a massive object will typically form a disk-like structure around the object. Artists' impressions such as the accompanying representation of a black hole with corona commonly depict the black hole as if it were a flat-space body hiding the part of the disk just behind it, but in reality gravitational lensing would greatly distort the image of the accretion disk.[160]
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+
146
+ Within such a disk, friction would cause angular momentum to be transported outward, allowing matter to fall further inward, thus releasing potential energy and increasing the temperature of the gas.[161]
147
+
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+ When the accreting object is a neutron star or a black hole, the gas in the inner accretion disk orbits at very high speeds because of its proximity to the compact object. The resulting friction is so significant that it heats the inner disk to temperatures at which it emits vast amounts of electromagnetic radiation (mainly X-rays). These bright X-ray sources may be detected by telescopes. This process of accretion is one of the most efficient energy-producing processes known; up to 40% of the rest mass of the accreted material can be emitted as radiation.[161] (In nuclear fusion only about 0.7% of the rest mass will be emitted as energy.) In many cases, accretion disks are accompanied by relativistic jets that are emitted along the poles, which carry away much of the energy. The mechanism for the creation of these jets is currently not well understood, in part due to insufficient data.[162]
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+
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+ As such, many of the universe's more energetic phenomena have been attributed to the accretion of matter on black holes. In particular, active galactic nuclei and quasars are believed to be the accretion disks of supermassive black holes.[163] Similarly, X-ray binaries are generally accepted to be binary star systems in which one of the two stars is a compact object accreting matter from its companion.[163] It has also been suggested that some ultraluminous X-ray sources may be the accretion disks of intermediate-mass black holes.[164]
151
+
152
+ In November 2011 the first direct observation of a quasar accretion disk around a supermassive black hole was reported.[165][166]
153
+
154
+ X-ray binaries are binary star systems that emit a majority of their radiation in the X-ray part of the spectrum. These X-ray emissions are generally thought to result when one of the stars (compact object) accretes matter from another (regular) star. The presence of an ordinary star in such a system provides an opportunity for studying the central object and to determine if it might be a black hole.[163]
155
+
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+ If such a system emits signals that can be directly traced back to the compact object, it cannot be a black hole. The absence of such a signal does, however, not exclude the possibility that the compact object is a neutron star. By studying the companion star it is often possible to obtain the orbital parameters of the system and to obtain an estimate for the mass of the compact object. If this is much larger than the Tolman–Oppenheimer–Volkoff limit (the maximum mass a star can have without collapsing) then the object cannot be a neutron star and is generally expected to be a black hole.[163]
157
+
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+ The first strong candidate for a black hole, Cygnus X-1, was discovered in this way by Charles Thomas Bolton,[167] Louise Webster and Paul Murdin[168] in 1972.[169][170] Some doubt, however, remained due to the uncertainties that result from the companion star being much heavier than the candidate black hole. Currently, better candidates for black holes are found in a class of X-ray binaries called soft X-ray transients. In this class of system, the companion star is of relatively low mass allowing for more accurate estimates of the black hole mass. Moreover, these systems actively emit X-rays for only several months once every 10–50 years. During the period of low X-ray emission (called quiescence), the accretion disk is extremely faint allowing detailed observation of the companion star during this period. One of the best such candidates is V404 Cygni.[163]
159
+
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+ The X-ray emissions from accretion disks sometimes flicker at certain frequencies. These signals are called quasi-periodic oscillations and are thought to be caused by material moving along the inner edge of the accretion disk (the innermost stable circular orbit). As such their frequency is linked to the mass of the compact object. They can thus be used as an alternative way to determine the mass of candidate black holes.[171]
161
+
162
+ Astronomers use the term "active galaxy" to describe galaxies with unusual characteristics, such as unusual spectral line emission and very strong radio emission. Theoretical and observational studies have shown that the activity in these active galactic nuclei (AGN) may be explained by the presence of supermassive black holes, which can be millions of times more massive than stellar ones. The models of these AGN consist of a central black hole that may be millions or billions of times more massive than the Sun; a disk of gas and dust called an accretion disk; and two jets perpendicular to the accretion disk.[172][173]
163
+
164
+ Although supermassive black holes are expected to be found in most AGN, only some galaxies' nuclei have been more carefully studied in attempts to both identify and measure the actual masses of the central supermassive black hole candidates. Some of the most notable galaxies with supermassive black hole candidates include the Andromeda Galaxy, M32, M87, NGC 3115, NGC 3377, NGC 4258, NGC 4889, NGC 1277, OJ 287, APM 08279+5255 and the Sombrero Galaxy.[175]
165
+
166
+ It is now widely accepted that the center of nearly every galaxy, not just active ones, contains a supermassive black hole.[176] The close observational correlation between the mass of this hole and the velocity dispersion of the host galaxy's bulge, known as the M-sigma relation, strongly suggests a connection between the formation of the black hole and the galaxy itself.[177]
167
+
168
+ Another way the black hole nature of an object may be tested in the future is through observation of effects caused by a strong gravitational field in their vicinity. One such effect is gravitational lensing: The deformation of spacetime around a massive object causes light rays to be deflected much as light passing through an optic lens. Observations have been made of weak gravitational lensing, in which light rays are deflected by only a few arcseconds. However, it has never been directly observed for a black hole.[179] One possibility for observing gravitational lensing by a black hole would be to observe stars in orbit around the black hole. There are several candidates for such an observation in orbit around Sagittarius A*.[179]
169
+
170
+ The evidence for stellar black holes strongly relies on the existence of an upper limit for the mass of a neutron star. The size of this limit heavily depends on the assumptions made about the properties of dense matter. New exotic phases of matter could push up this bound.[163] A phase of free quarks at high density might allow the existence of dense quark stars,[180] and some supersymmetric models predict the existence of Q stars.[181] Some extensions of the standard model posit the existence of preons as fundamental building blocks of quarks and leptons, which could hypothetically form preon stars.[182] These hypothetical models could potentially explain a number of observations of stellar black hole candidates. However, it can be shown from arguments in general relativity that any such object will have a maximum mass.[163]
171
+
172
+ Since the average density of a black hole inside its Schwarzschild radius is inversely proportional to the square of its mass, supermassive black holes are much less dense than stellar black holes (the average density of a 108 M☉ black hole is comparable to that of water).[163] Consequently, the physics of matter forming a supermassive black hole is much better understood and the possible alternative explanations for supermassive black hole observations are much more mundane. For example, a supermassive black hole could be modelled by a large cluster of very dark objects. However, such alternatives are typically not stable enough to explain the supermassive black hole candidates.[163]
173
+
174
+ The evidence for the existence of stellar and supermassive black holes implies that in order for black holes to not form, general relativity must fail as a theory of gravity, perhaps due to the onset of quantum mechanical corrections. A much anticipated feature of a theory of quantum gravity is that it will not feature singularities or event horizons and thus black holes would not be real artifacts.[183] For example, in the fuzzball model based on string theory, the individual states of a black hole solution do not generally have an event horizon or singularity, but for a classical/semi-classical observer the statistical average of such states appears just as an ordinary black hole as deduced from general relativity.[184]
175
+
176
+ A few theoretical objects have been conjectured to match observations of astronomical black hole candidates identically or near-identically, but which function via a different mechanism. These include the gravastar, the black star,[185] and the dark-energy star.[186]
177
+
178
+ In 1971, Hawking showed under general conditions[Note 5] that the total area of the event horizons of any collection of classical black holes can never decrease, even if they collide and merge.[187] This result, now known as the second law of black hole mechanics, is remarkably similar to the second law of thermodynamics, which states that the total entropy of an isolated system can never decrease. As with classical objects at absolute zero temperature, it was assumed that black holes had zero entropy. If this were the case, the second law of thermodynamics would be violated by entropy-laden matter entering a black hole, resulting in a decrease of the total entropy of the universe. Therefore, Bekenstein proposed that a black hole should have an entropy, and that it should be proportional to its horizon area.[188]
179
+
180
+ The link with the laws of thermodynamics was further strengthened by Hawking's discovery that quantum field theory predicts that a black hole radiates blackbody radiation at a constant temperature. This seemingly causes a violation of the second law of black hole mechanics, since the radiation will carry away energy from the black hole causing it to shrink. The radiation, however also carries away entropy, and it can be proven under general assumptions that the sum of the entropy of the matter surrounding a black hole and one quarter of the area of the horizon as measured in Planck units is in fact always increasing. This allows the formulation of the first law of black hole mechanics as an analogue of the first law of thermodynamics, with the mass acting as energy, the surface gravity as temperature and the area as entropy.[188]
181
+
182
+ One puzzling feature is that the entropy of a black hole scales with its area rather than with its volume, since entropy is normally an extensive quantity that scales linearly with the volume of the system. This odd property led Gerard 't Hooft and Leonard Susskind to propose the holographic principle, which suggests that anything that happens in a volume of spacetime can be described by data on the boundary of that volume.[189]
183
+
184
+ Although general relativity can be used to perform a semi-classical calculation of black hole entropy, this situation is theoretically unsatisfying. In statistical mechanics, entropy is understood as counting the number of microscopic configurations of a system that have the same macroscopic qualities (such as mass, charge, pressure, etc.). Without a satisfactory theory of quantum gravity, one cannot perform such a computation for black holes. Some progress has been made in various approaches to quantum gravity. In 1995, Andrew Strominger and Cumrun Vafa showed that counting the microstates of a specific supersymmetric black hole in string theory reproduced the Bekenstein–Hawking entropy.[190] Since then, similar results have been reported for different black holes both in string theory and in other approaches to quantum gravity like loop quantum gravity.[191]
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+
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+ Because a black hole has only a few internal parameters, most of the information about the matter that went into forming the black hole is lost. Regardless of the type of matter which goes into a black hole, it appears that only information concerning the total mass, charge, and angular momentum are conserved. As long as black holes were thought to persist forever this information loss is not that problematic, as the information can be thought of as existing inside the black hole, inaccessible from the outside, but represented on the event horizon in accordance with the holographic principle. However, black holes slowly evaporate by emitting Hawking radiation. This radiation does not appear to carry any additional information about the matter that formed the black hole, meaning that this information appears to be gone forever.[192]
187
+
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+ The question whether information is truly lost in black holes (the black hole information paradox) has divided the theoretical physics community (see Thorne–Hawking–Preskill bet). In quantum mechanics, loss of information corresponds to the violation of a property called unitarity, and it has been argued that loss of unitarity would also imply violation of conservation of energy,[193] though this has also been disputed.[194] Over recent years evidence has been building that indeed information and unitarity are preserved in a full quantum gravitational treatment of the problem.[195]
189
+
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+ One attempt to resolve the black hole information paradox is known as black hole complementarity. In 2012, the "firewall paradox" was introduced with the goal of demonstrating that black hole complementarity fails to solve the information paradox. According to quantum field theory in curved spacetime, a single emission of Hawking radiation involves two mutually entangled particles. The outgoing particle escapes and is emitted as a quantum of Hawking radiation; the infalling particle is swallowed by the black hole. Assume a black hole formed a finite time in the past and will fully evaporate away in some finite time in the future. Then, it will emit only a finite amount of information encoded within its Hawking radiation. According to research by physicists like Don Page[196][197] and Leonard Susskind, there will eventually be a time by which an outgoing particle must be entangled with all the Hawking radiation the black hole has previously emitted. This seemingly creates a paradox: a principle called "monogamy of entanglement" requires that, like any quantum system, the outgoing particle cannot be fully entangled with two other systems at the same time; yet here the outgoing particle appears to be entangled both with the infalling particle and, independently, with past Hawking radiation.[198] In order to resolve this contradiction, physicists may eventually be forced to give up one of three time-tested principles: Einstein's equivalence principle, unitarity, or local quantum field theory. One possible solution, which violates the equivalence principle, is that a "firewall" destroys incoming particles at the event horizon.[199] In general, which if any of these assumptions should be abandoned remains a topic of debate.[194]
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+
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1
+
2
+
3
+ A black hole is a region of spacetime where gravity is so strong that nothing—no particles or even electromagnetic radiation such as light—can escape from it.[6] The theory of general relativity predicts that a sufficiently compact mass can deform spacetime to form a black hole.[7][8] The boundary of the region from which no escape is possible is called the event horizon. Although the event horizon has an enormous effect on the fate and circumstances of an object crossing it, according to general relativity it has no locally detectable features.[9] In many ways, a black hole acts like an ideal black body, as it reflects no light.[10][11] Moreover, quantum field theory in curved spacetime predicts that event horizons emit Hawking radiation, with the same spectrum as a black body of a temperature inversely proportional to its mass. This temperature is on the order of billionths of a kelvin for black holes of stellar mass, making it essentially impossible to observe.
4
+
5
+ Objects whose gravitational fields are too strong for light to escape were first considered in the 18th century by John Michell and Pierre-Simon Laplace.[12] The first modern solution of general relativity that would characterize a black hole was found by Karl Schwarzschild in 1916, although its interpretation as a region of space from which nothing can escape was first published by David Finkelstein in 1958. Black holes were long considered a mathematical curiosity; it was not until the 1960s that theoretical work showed they were a generic prediction of general relativity. The discovery of neutron stars by Jocelyn Bell Burnell in 1967 sparked interest in gravitationally collapsed compact objects as a possible astrophysical reality.
6
+
7
+ Black holes of stellar mass are expected to form when very massive stars collapse at the end of their life cycle. After a black hole has formed, it can continue to grow by absorbing mass from its surroundings. By absorbing other stars and merging with other black holes, supermassive black holes of millions of solar masses (M☉) may form. There is consensus that supermassive black holes exist in the centers of most galaxies.
8
+
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+ The presence of a black hole can be inferred through its interaction with other matter and with electromagnetic radiation such as visible light. Matter that falls onto a black hole can form an external accretion disk heated by friction, forming quasars, some of the brightest objects in the universe. Stars passing too close to a supermassive black hole can be shred into streamers that shine very brightly before being "swallowed."[13] If there are other stars orbiting a black hole, their orbits can be used to determine the black hole's mass and location. Such observations can be used to exclude possible alternatives such as neutron stars. In this way, astronomers have identified numerous stellar black hole candidates in binary systems, and established that the radio source known as Sagittarius A*, at the core of the Milky Way galaxy, contains a supermassive black hole of about 4.3 million solar masses.
10
+
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+ On 11 February 2016, the LIGO Scientific Collaboration and the Virgo collaboration announced the first direct detection of gravitational waves, which also represented the first observation of a black hole merger.[14] As of December 2018[update], eleven gravitational wave events have been observed that originated from ten merging black holes (along with one binary neutron star merger).[15][16] On 10 April 2019, the first direct image of a black hole and its vicinity was published, following observations made by the Event Horizon Telescope in 2017 of the supermassive black hole in Messier 87's galactic centre.[3][17][18]
12
+
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+ The idea of a body so massive that even light could not escape was briefly proposed by astronomical pioneer and English clergyman John Michell in a letter published in November 1784. Michell's simplistic calculations assumed such a body might have the same density as the Sun, and concluded that such a body would form when a star's diameter exceeds the Sun's by a factor of 500, and the surface escape velocity exceeds the usual speed of light. Michell correctly noted that such supermassive but non-radiating bodies might be detectable through their gravitational effects on nearby visible bodies.[20][12][21] Scholars of the time were initially excited by the proposal that giant but invisible stars might be hiding in plain view, but enthusiasm dampened when the wavelike nature of light became apparent in the early nineteenth century.[22]
14
+
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+ If light were a wave rather than a "corpuscle", it is unclear what, if any, influence gravity would have on escaping light waves.[12][21] Modern physics discredits Michell's notion of a light ray shooting directly from the surface of a supermassive star, being slowed down by the star's gravity, stopping, and then free-falling back to the star's surface.[23]
16
+
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+ In 1915, Albert Einstein developed his theory of general relativity, having earlier shown that gravity does influence light's motion. Only a few months later, Karl Schwarzschild found a solution to the Einstein field equations, which describes the gravitational field of a point mass and a spherical mass.[24] A few months after Schwarzschild, Johannes Droste, a student of Hendrik Lorentz, independently gave the same solution for the point mass and wrote more extensively about its properties.[25][26] This solution had a peculiar behaviour at what is now called the Schwarzschild radius, where it became singular, meaning that some of the terms in the Einstein equations became infinite. The nature of this surface was not quite understood at the time. In 1924, Arthur Eddington showed that the singularity disappeared after a change of coordinates (see Eddington–Finkelstein coordinates), although it took until 1933 for Georges Lemaître to realize that this meant the singularity at the Schwarzschild radius was a non-physical coordinate singularity.[27] Arthur Eddington did however comment on the possibility of a star with mass compressed to the Schwarzschild radius in a 1926 book, noting that Einstein's theory allows us to rule out overly large densities for visible stars like Betelgeuse because "a star of 250 million km radius could not possibly have so high a density as the sun. Firstly, the force of gravitation would be so great that light would be unable to escape from it, the rays falling back to the star like a stone to the earth. Secondly, the red shift of the spectral lines would be so great that the spectrum would be shifted out of existence. Thirdly, the mass would produce so much curvature of the space-time metric that space would close up around the star, leaving us outside (i.e., nowhere)."[28][29]
18
+
19
+ In 1931, Subrahmanyan Chandrasekhar calculated, using special relativity, that a non-rotating body of electron-degenerate matter above a certain limiting mass (now called the Chandrasekhar limit at 1.4 M☉) has no stable solutions.[30] His arguments were opposed by many of his contemporaries like Eddington and Lev Landau, who argued that some yet unknown mechanism would stop the collapse.[31] They were partly correct: a white dwarf slightly more massive than the Chandrasekhar limit will collapse into a neutron star,[32] which is itself stable. But in 1939, Robert Oppenheimer and others predicted that neutron stars above another limit (the Tolman–Oppenheimer–Volkoff limit) would collapse further for the reasons presented by Chandrasekhar, and concluded that no law of physics was likely to intervene and stop at least some stars from collapsing to black holes.[33] Their original calculations, based on the Pauli exclusion principle, gave it as 0.7 M☉; subsequent consideration of strong force-mediated neutron-neutron repulsion raised the estimate to approximately 1.5 M☉ to 3.0 M☉.[34] Observations of the neutron star merger GW170817, which is thought to have generated a black hole shortly afterward, have refined the TOV limit estimate to ~2.17 M☉.[35][36][37][38][39]
20
+
21
+ Oppenheimer and his co-authors interpreted the singularity at the boundary of the Schwarzschild radius as indicating that this was the boundary of a bubble in which time stopped. This is a valid point of view for external observers, but not for infalling observers. Because of this property, the collapsed stars were called "frozen stars", because an outside observer would see the surface of the star frozen in time at the instant where its collapse takes it to the Schwarzschild radius.[40]
22
+
23
+ In 1958, David Finkelstein identified the Schwarzschild surface as an event horizon, "a perfect unidirectional membrane: causal influences can cross it in only one direction".[41] This did not strictly contradict Oppenheimer's results, but extended them to include the point of view of infalling observers. Finkelstein's solution extended the Schwarzschild solution for the future of observers falling into a black hole. A complete extension had already been found by Martin Kruskal, who was urged to publish it.[42]
24
+
25
+ These results came at the beginning of the golden age of general relativity, which was marked by general relativity and black holes becoming mainstream subjects of research. This process was helped by the discovery of pulsars by Jocelyn Bell Burnell in 1967,[43][44] which, by 1969, were shown to be rapidly rotating neutron stars.[45] Until that time, neutron stars, like black holes, were regarded as just theoretical curiosities; but the discovery of pulsars showed their physical relevance and spurred a further interest in all types of compact objects that might be formed by gravitational collapse.[citation needed]
26
+
27
+ In this period more general black hole solutions were found. In 1963, Roy Kerr found the exact solution for a rotating black hole. Two years later, Ezra Newman found the axisymmetric solution for a black hole that is both rotating and electrically charged.[46] Through the work of Werner Israel,[47] Brandon Carter,[48][49] and David Robinson[50] the no-hair theorem emerged, stating that a stationary black hole solution is completely described by the three parameters of the Kerr–Newman metric: mass, angular momentum, and electric charge.[51]
28
+
29
+ At first, it was suspected that the strange features of the black hole solutions were pathological artifacts from the symmetry conditions imposed, and that the singularities would not appear in generic situations. This view was held in particular by Vladimir Belinsky, Isaak Khalatnikov, and Evgeny Lifshitz, who tried to prove that no singularities appear in generic solutions. However, in the late 1960s Roger Penrose[52] and Stephen Hawking used global techniques to prove that singularities appear generically.[53]
30
+
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+ Work by James Bardeen, Jacob Bekenstein, Carter, and Hawking in the early 1970s led to the formulation of black hole thermodynamics.[54] These laws describe the behaviour of a black hole in close analogy to the laws of thermodynamics by relating mass to energy, area to entropy, and surface gravity to temperature. The analogy was completed when Hawking, in 1974, showed that quantum field theory implies that black holes should radiate like a black body with a temperature proportional to the surface gravity of the black hole, predicting the effect now known as Hawking radiation.[55]
32
+
33
+ John Michell used the term "dark star",[56] and in the early 20th century, physicists used the term "gravitationally collapsed object". Science writer Marcia Bartusiak traces the term "black hole" to physicist Robert H. Dicke, who in the early 1960s reportedly compared the phenomenon to the Black Hole of Calcutta, notorious as a prison where people entered but never left alive.[57]
34
+
35
+ The term "black hole" was used in print by Life and Science News magazines in 1963,[57] and by science journalist Ann Ewing in her article "'Black Holes' in Space", dated 18 January 1964, which was a report on a meeting of the American Association for the Advancement of Science held in Cleveland, Ohio.[58][59]
36
+
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+ In December 1967, a student reportedly suggested the phrase "black hole" at a lecture by John Wheeler;[58] Wheeler adopted the term for its brevity and "advertising value", and it quickly caught on,[60] leading some to credit Wheeler with coining the phrase.[61]
38
+
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+ The no-hair conjecture postulates that, once it achieves a stable condition after formation, a black hole has only three independent physical properties: mass, charge, and angular momentum; the black hole is otherwise featureless. If the conjecture is true, any two black holes that share the same values for these properties, or parameters, are indistinguishable from one another. The degree to which the conjecture is true for real black holes under the laws of modern physics, is currently an unsolved problem.[51]
40
+
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+ These properties are special because they are visible from outside a black hole. For example, a charged black hole repels other like charges just like any other charged object. Similarly, the total mass inside a sphere containing a black hole can be found by using the gravitational analog of Gauss's law, the ADM mass, far away from the black hole.[62][clarification needed] Likewise, the angular momentum can be measured from far away using frame dragging by the gravitomagnetic field.[clarification needed]
42
+
43
+ When an object falls into a black hole, any information about the shape of the object or distribution of charge on it is evenly distributed along the horizon of the black hole, and is lost to outside observers. The behavior of the horizon in this situation is a dissipative system that is closely analogous to that of a conductive stretchy membrane with friction and electrical resistance—the membrane paradigm.[63] This is different from other field theories such as electromagnetism, which do not have any friction or resistivity at the microscopic level, because they are time-reversible. Because a black hole eventually achieves a stable state with only three parameters, there is no way to avoid losing information about the initial conditions: the gravitational and electric fields of a black hole give very little information about what went in. The information that is lost includes every quantity that cannot be measured far away from the black hole horizon, including approximately conserved quantum numbers such as the total baryon number and lepton number. This behavior is so puzzling that it has been called the black hole information loss paradox.[64][65]
44
+
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+ The simplest static black holes have mass but neither electric charge nor angular momentum. These black holes are often referred to as Schwarzschild black holes after Karl Schwarzschild who discovered this solution in 1916.[24] According to Birkhoff's theorem, it is the only vacuum solution that is spherically symmetric.[66] This means there is no observable difference at a distance between the gravitational field of such a black hole and that of any other spherical object of the same mass. The popular notion of a black hole "sucking in everything" in its surroundings is therefore correct only near a black hole's horizon; far away, the external gravitational field is identical to that of any other body of the same mass.[67]
46
+
47
+ Solutions describing more general black holes also exist. Non-rotating charged black holes are described by the Reissner–Nordström metric, while the Kerr metric describes a non-charged rotating black hole. The most general stationary black hole solution known is the Kerr–Newman metric, which describes a black hole with both charge and angular momentum.[68]
48
+
49
+ While the mass of a black hole can take any positive value, the charge and angular momentum are constrained by the mass. In Planck units, the total electric charge Q and the total angular momentum J are expected to satisfy
50
+
51
+ for a black hole of mass M. Black holes with the minimum possible mass satisfying this inequality are called extremal. Solutions of Einstein's equations that violate this inequality exist, but they do not possess an event horizon. These solutions have so-called naked singularities that can be observed from the outside, and hence are deemed unphysical. The cosmic censorship hypothesis rules out the formation of such singularities, when they are created through the gravitational collapse of realistic matter.[7] This is supported by numerical simulations.[69]
52
+
53
+ Due to the relatively large strength of the electromagnetic force, black holes forming from the collapse of stars are expected to retain the nearly neutral charge of the star. Rotation, however, is expected to be a universal feature of compact astrophysical objects. The black-hole candidate binary X-ray source GRS 1915+105[70] appears to have an angular momentum near the maximum allowed value. That uncharged limit is[71]
54
+
55
+ allowing definition of a dimensionless spin parameter such that[71]
56
+
57
+ Black holes are commonly classified according to their mass, independent of angular momentum, J. The size of a black hole, as determined by the radius of the event horizon, or Schwarzschild radius, is proportional to the mass, M, through
58
+
59
+ where rs is the Schwarzschild radius and MSun is the mass of the Sun.[73] For a black hole with nonzero spin and/or electric charge, the radius is smaller,[Note 2] until an extremal black hole could have an event horizon close to[74]
60
+
61
+ The defining feature of a black hole is the appearance of an event horizon—a boundary in spacetime through which matter and light can pass only inward towards the mass of the black hole. Nothing, not even light, can escape from inside the event horizon. The event horizon is referred to as such because if an event occurs within the boundary, information from that event cannot reach an outside observer, making it impossible to determine whether such an event occurred.[76]
62
+
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+ As predicted by general relativity, the presence of a mass deforms spacetime in such a way that the paths taken by particles bend towards the mass.[77] At the event horizon of a black hole, this deformation becomes so strong that there are no paths that lead away from the black hole.[78]
64
+
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+ To a distant observer, clocks near a black hole would appear to tick more slowly than those further away from the black hole.[79] Due to this effect, known as gravitational time dilation, an object falling into a black hole appears to slow as it approaches the event horizon, taking an infinite time to reach it.[80] At the same time, all processes on this object slow down, from the view point of a fixed outside observer, causing any light emitted by the object to appear redder and dimmer, an effect known as gravitational redshift.[81] Eventually, the falling object fades away until it can no longer be seen. Typically this process happens very rapidly with an object disappearing from view within less than a second.[82]
66
+
67
+ On the other hand, indestructible observers falling into a black hole do not notice any of these effects as they cross the event horizon. According to their own clocks, which appear to them to tick normally, they cross the event horizon after a finite time without noting any singular behaviour; in classical general relativity, it is impossible to determine the location of the event horizon from local observations, due to Einstein's equivalence principle.[83][84]
68
+
69
+ The topology of the event horizon of a black hole at equilibrium is always spherical.[Note 4][87] For non-rotating (static) black holes the geometry of the event horizon is precisely spherical, while for rotating black holes the event horizon is oblate.[88][89][90]
70
+
71
+ At the center of a black hole, as described by general relativity, may lie a gravitational singularity, a region where the spacetime curvature becomes infinite.[91] For a non-rotating black hole, this region takes the shape of a single point and for a rotating black hole, it is smeared out to form a ring singularity that lies in the plane of rotation.[92] In both cases, the singular region has zero volume. It can also be shown that the singular region contains all the mass of the black hole solution.[93] The singular region can thus be thought of as having infinite density.[94]
72
+
73
+ Observers falling into a Schwarzschild black hole (i.e., non-rotating and not charged) cannot avoid being carried into the singularity once they cross the event horizon. They can prolong the experience by accelerating away to slow their descent, but only up to a limit.[95] When they reach the singularity, they are crushed to infinite density and their mass is added to the total of the black hole. Before that happens, they will have been torn apart by the growing tidal forces in a process sometimes referred to as spaghettification or the "noodle effect".[96]
74
+
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+ In the case of a charged (Reissner–Nordström) or rotating (Kerr) black hole, it is possible to avoid the singularity. Extending these solutions as far as possible reveals the hypothetical possibility of exiting the black hole into a different spacetime with the black hole acting as a wormhole.[97] The possibility of traveling to another universe is, however, only theoretical since any perturbation would destroy this possibility.[98] It also appears to be possible to follow closed timelike curves (returning to one's own past) around the Kerr singularity, which leads to problems with causality like the grandfather paradox.[99] It is expected that none of these peculiar effects would survive in a proper quantum treatment of rotating and charged black holes.[100]
76
+
77
+ The appearance of singularities in general relativity is commonly perceived as signaling the breakdown of the theory.[101] This breakdown, however, is expected; it occurs in a situation where quantum effects should describe these actions, due to the extremely high density and therefore particle interactions. To date, it has not been possible to combine quantum and gravitational effects into a single theory, although there exist attempts to formulate such a theory of quantum gravity. It is generally expected that such a theory will not feature any singularities.[102][103]
78
+
79
+ The photon sphere is a spherical boundary of zero thickness in which photons that move on tangents to that sphere would be trapped in a circular orbit about the black hole. For non-rotating black holes, the photon sphere has a radius 1.5 times the Schwarzschild radius. Their orbits would be dynamically unstable, hence any small perturbation, such as a particle of infalling matter, would cause an instability that would grow over time, either setting the photon on an outward trajectory causing it to escape the black hole, or on an inward spiral where it would eventually cross the event horizon.[104]
80
+
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+ While light can still escape from the photon sphere, any light that crosses the photon sphere on an inbound trajectory will be captured by the black hole. Hence any light that reaches an outside observer from the photon sphere must have been emitted by objects between the photon sphere and the event horizon.[104]
82
+
83
+ Rotating black holes are surrounded by a region of spacetime in which it is impossible to stand still, called the ergosphere. This is the result of a process known as frame-dragging; general relativity predicts that any rotating mass will tend to slightly "drag" along the spacetime immediately surrounding it. Any object near the rotating mass will tend to start moving in the direction of rotation. For a rotating black hole, this effect is so strong near the event horizon that an object would have to move faster than the speed of light in the opposite direction to just stand still.[106]
84
+
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+ The ergosphere of a black hole is a volume whose inner boundary is the black hole's event horizon and an outer boundary called the ergosurface, which coincides with the event horizon at the poles but noticeably wider around the equator.[105]
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+
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+ Objects and radiation can escape normally from the ergosphere. Through the Penrose process, objects can emerge from the ergosphere with more energy than they entered. This energy is taken from the rotational energy of the black hole causing the latter to slow.[107] A variation of the Penrose process in the presence of strong magnetic fields, the Blandford–Znajek process is considered a likely mechanism for the enormous luminosity and relativistic jets of quasars and other active galactic nuclei.
88
+
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+ In Newtonian gravity, test particles can stably orbit at arbitrary distances from a central object. In general relativity, however, there exists an innermost stable circular orbit (often called the ISCO), inside of which, any infinitesimal perturbations to a circular orbit will lead to inspiral into the black hole.[108] The location of the ISCO depends on the spin of the black hole, in the case of a Schwarzschild black hole (spin zero) is:
90
+
91
+ and decreases with increasing black hole spin for particles orbiting in the same direction as the spin.[109]
92
+
93
+ Given the bizarre character of black holes, it was long questioned whether such objects could actually exist in nature or whether they were merely pathological solutions to Einstein's equations. Einstein himself wrongly thought black holes would not form, because he held that the angular momentum of collapsing particles would stabilize their motion at some radius.[110] This led the general relativity community to dismiss all results to the contrary for many years. However, a minority of relativists continued to contend that black holes were physical objects,[111] and by the end of the 1960s, they had persuaded the majority of researchers in the field that there is no obstacle to the formation of an event horizon.[citation needed]
94
+
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+ Penrose demonstrated that once an event horizon forms, general relativity without quantum mechanics requires that a singularity will form within.[52] Shortly afterwards, Hawking showed that many cosmological solutions that describe the Big Bang have singularities without scalar fields or other exotic matter (see "Penrose–Hawking singularity theorems").[clarification needed] The Kerr solution, the no-hair theorem, and the laws of black hole thermodynamics showed that the physical properties of black holes were simple and comprehensible, making them respectable subjects for research.[112] Conventional black holes are formed by gravitational collapse of heavy objects such as stars, but they can also in theory be formed by other processes.[113][114]
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+
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+ Gravitational collapse occurs when an object's internal pressure is insufficient to resist the object's own gravity. For stars this usually occurs either because a star has too little "fuel" left to maintain its temperature through stellar nucleosynthesis, or because a star that would have been stable receives extra matter in a way that does not raise its core temperature. In either case the star's temperature is no longer high enough to prevent it from collapsing under its own weight.[115]
98
+ The collapse may be stopped by the degeneracy pressure of the star's constituents, allowing the condensation of matter into an exotic denser state. The result is one of the various types of compact star. Which type forms depends on the mass of the remnant of the original star left after the outer layers have been blown away. Such explosions and pulsations lead to planetary nebula.[116] This mass can be substantially less than the original star. Remnants exceeding 5 M☉ are produced by stars that were over 20 M☉ before the collapse.[115]
99
+
100
+ If the mass of the remnant exceeds about 3–4 M☉ (the Tolman–Oppenheimer–Volkoff limit[33]), either because the original star was very heavy or because the remnant collected additional mass through accretion of matter, even the degeneracy pressure of neutrons is insufficient to stop the collapse. No known mechanism (except possibly quark degeneracy pressure, see quark star) is powerful enough to stop the implosion and the object will inevitably collapse to form a black hole.[115]
101
+
102
+ The gravitational collapse of heavy stars is assumed to be responsible for the formation of stellar mass black holes. Star formation in the early universe may have resulted in very massive stars, which upon their collapse would have produced black holes of up to 103 M☉. These black holes could be the seeds of the supermassive black holes found in the centers of most galaxies.[118] It has further been suggested that supermassive black holes with typical masses of ~105 M☉ could have formed from the direct collapse of gas clouds in the young universe.[113] Some candidates for such objects have been found in observations of the young universe.[113]
103
+
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+ While most of the energy released during gravitational collapse is emitted very quickly, an outside observer does not actually see the end of this process. Even though the collapse takes a finite amount of time from the reference frame of infalling matter, a distant observer would see the infalling material slow and halt just above the event horizon, due to gravitational time dilation. Light from the collapsing material takes longer and longer to reach the observer, with the light emitted just before the event horizon forms delayed an infinite amount of time. Thus the external observer never sees the formation of the event horizon; instead, the collapsing material seems to become dimmer and increasingly red-shifted, eventually fading away.[119]
105
+
106
+ Gravitational collapse requires great density. In the current epoch of the universe these high densities are found only in stars, but in the early universe shortly after the Big Bang densities were much greater, possibly allowing for the creation of black holes. High density alone is not enough to allow black hole formation since a uniform mass distribution will not allow the mass to bunch up. In order for primordial black holes to have formed in such a dense medium, there must have been initial density perturbations that could then grow under their own gravity. Different models for the early universe vary widely in their predictions of the scale of these fluctuations. Various models predict the creation of primordial black holes ranging in size from a Planck mass to hundreds of thousands of solar masses.[114]
107
+
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+ Despite the early universe being extremely dense—far denser than is usually required to form a black hole—it did not re-collapse into a black hole during the Big Bang. Models for gravitational collapse of objects of relatively constant size, such as stars, do not necessarily apply in the same way to rapidly expanding space such as the Big Bang.[120]
109
+
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+ Gravitational collapse is not the only process that could create black holes. In principle, black holes could be formed in high-energy collisions that achieve sufficient density. As of 2002, no such events have been detected, either directly or indirectly as a deficiency of the mass balance in particle accelerator experiments.[121] This suggests that there must be a lower limit for the mass of black holes. Theoretically, this boundary is expected to lie around the Planck mass (mP=√ħ c/G ≈ 1.2×1019 GeV/c2 ≈ 2.2×10−8 kg), where quantum effects are expected to invalidate the predictions of general relativity.[122] This would put the creation of black holes firmly out of reach of any high-energy process occurring on or near the Earth. However, certain developments in quantum gravity suggest that the minimum black hole mass could be much lower: some braneworld scenarios for example put the boundary as low as 1 TeV/c2.[123] This would make it conceivable for micro black holes to be created in the high-energy collisions that occur when cosmic rays hit the Earth's atmosphere, or possibly in the Large Hadron Collider at CERN. These theories are very speculative, and the creation of black holes in these processes is deemed unlikely by many specialists.[124] Even if micro black holes could be formed, it is expected that they would evaporate in about 10−25 seconds, posing no threat to the Earth.[125]
111
+
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+ Once a black hole has formed, it can continue to grow by absorbing additional matter. Any black hole will continually absorb gas and interstellar dust from its surroundings. This is the primary process through which supermassive black holes seem to have grown.[118] A similar process has been suggested for the formation of intermediate-mass black holes found in globular clusters.[126] Black holes can also merge with other objects such as stars or even other black holes. This is thought to have been important, especially in the early growth of supermassive black holes, which could have formed from the aggregation of many smaller objects.[118] The process has also been proposed as the origin of some intermediate-mass black holes.[127][128]
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+
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+ In 1974, Hawking predicted that black holes are not entirely black but emit small amounts of thermal radiation at a temperature ℏ c3/(8 π G M kB);[55] this effect has become known as Hawking radiation. By applying quantum field theory to a static black hole background, he determined that a black hole should emit particles that display a perfect black body spectrum. Since Hawking's publication, many others have verified the result through various approaches.[129] If Hawking's theory of black hole radiation is correct, then black holes are expected to shrink and evaporate over time as they lose mass by the emission of photons and other particles.[55] The temperature of this thermal spectrum (Hawking temperature) is proportional to the surface gravity of the black hole, which, for a Schwarzschild black hole, is inversely proportional to the mass. Hence, large black holes emit less radiation than small black holes.[130]
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+ A stellar black hole of 1 M☉ has a Hawking temperature of 62 nanokelvins.[131] This is far less than the 2.7 K temperature of the cosmic microwave background radiation. Stellar-mass or larger black holes receive more mass from the cosmic microwave background than they emit through Hawking radiation and thus will grow instead of shrinking.[132] To have a Hawking temperature larger than 2.7 K (and be able to evaporate), a black hole would need a mass less than the Moon. Such a black hole would have a diameter of less than a tenth of a millimeter.[133]
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+
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+ If a black hole is very small, the radiation effects are expected to become very strong. A black hole with the mass of a car would have a diameter of about 10−24 m and take a nanosecond to evaporate, during which time it would briefly have a luminosity of more than 200 times that of the Sun. Lower-mass black holes are expected to evaporate even faster; for example, a black hole of mass 1 TeV/c2 would take less than 10−88 seconds to evaporate completely. For such a small black hole, quantum gravitation effects are expected to play an important role and could hypothetically make such a small black hole stable, although current developments in quantum gravity do not indicate this is the case.[134][135]
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+ The Hawking radiation for an astrophysical black hole is predicted to be very weak and would thus be exceedingly difficult to detect from Earth. A possible exception, however, is the burst of gamma rays emitted in the last stage of the evaporation of primordial black holes. Searches for such flashes have proven unsuccessful and provide stringent limits on the possibility of existence of low mass primordial black holes.[136] NASA's Fermi Gamma-ray Space Telescope launched in 2008 will continue the search for these flashes.[137]
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+ If black holes evaporate via Hawking radiation, a solar mass black hole will evaporate (beginning once the temperature of the cosmic microwave background drops below that of the black hole) over a period of 1064 years.[138] A supermassive black hole with a mass of 1011 (100 billion) M☉ will evaporate in around 2×10100 years.[139] Some monster black holes in the universe are predicted to continue to grow up to perhaps 1014 M☉ during the collapse of superclusters of galaxies. Even these would evaporate over a timescale of up to 10106 years.[138]
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+ By nature, black holes do not themselves emit any electromagnetic radiation other than the hypothetical Hawking radiation, so astrophysicists searching for black holes must generally rely on indirect observations. For example, a black hole's existence can sometimes be inferred by observing its gravitational influence upon its surroundings.[140]
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+
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+ On 10 April 2019 an image was released of a black hole, which is seen in magnified fashion because the light paths near the event horizon are highly bent. The dark shadow in the middle results from light paths absorbed by the black hole. The image is in false color, as the detected light halo in this image is not in the visible spectrum, but radio waves.
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+ The Event Horizon Telescope (EHT), is an active program that directly observes the immediate environment of the event horizon of black holes, such as the black hole at the centre of the Milky Way. In April 2017, EHT began observation of the black hole in the center of Messier 87.[141] "In all, eight radio observatories on six mountains and four continents observed the galaxy in Virgo on and off for 10 days in April 2017" to provide the data yielding the image two years later in April 2019.[142] After two years of data processing, EHT released the first direct image of a black hole, specifically the supermassive black hole that lies in the center of the aforementioned galaxy.[143][144] What is visible is not the black hole, which shows as black because of the loss of all light within this dark region, rather it is the gases at the edge of the event horizon, which are displayed as orange or red, that define the black hole.[145]
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+
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+ The brightening of this material in the 'bottom' half of the processed EHT image is thought to be caused by Doppler beaming, whereby material approaching the viewer at relativistic speeds is perceived as brighter than material moving away. In the case of a black hole this phenomenon implies that the visible material is rotating at relativistic speeds (>1,000 km/s), the only speeds at which it is possible to centrifugally balance the immense gravitational attraction of the singularity, and thereby remain in orbit above the event horizon. This configuration of bright material implies that the EHT observed M87* from a perspective catching the black hole's accretion disc nearly edge-on, as the whole system rotated clockwise.[146] However, the extreme gravitational lensing associated with black holes produces the illusion of a perspective that sees the accretion disc from above. In reality, most of the ring in the EHT image was created when the light emitted by the far side of the accretion disc bent around the black hole's gravity well and escaped such that most of the possible perspectives on M87* can see the entire disc, even that directly behind the "shadow".
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+ Prior to this, in 2015, the EHT detected magnetic fields just outside the event horizon of Sagittarius A*, and even discerned some of their properties. The field lines that pass through the accretion disc were found to be a complex mixture of ordered and tangled. The existence of magnetic fields had been predicted by theoretical studies of black holes.[147][148]
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+
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+ On 14 September 2015 the LIGO gravitational wave observatory made the first-ever successful direct observation of gravitational waves.[14][150] The signal was consistent with theoretical predictions for the gravitational waves produced by the merger of two black holes: one with about 36 solar masses, and the other around 29 solar masses.[14][151] This observation provides the most concrete evidence for the existence of black holes to date. For instance, the gravitational wave signal suggests that the separation of the two objects prior to the merger was just 350 km (or roughly four times the Schwarzschild radius corresponding to the inferred masses). The objects must therefore have been extremely compact, leaving black holes as the most plausible interpretation.[14]
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+
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+ More importantly, the signal observed by LIGO also included the start of the post-merger ringdown, the signal produced as the newly formed compact object settles down to a stationary state. Arguably, the ringdown is the most direct way of observing a black hole.[152] From the LIGO signal it is possible to extract the frequency and damping time of the dominant mode of the ringdown. From these it is possible to infer the mass and angular momentum of the final object, which match independent predictions from numerical simulations of the merger.[153] The frequency and decay time of the dominant mode are determined by the geometry of the photon sphere. Hence, observation of this mode confirms the presence of a photon sphere, however it cannot exclude possible exotic alternatives to black holes that are compact enough to have a photon sphere.[152]
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+
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+ The observation also provides the first observational evidence for the existence of stellar-mass black hole binaries. Furthermore, it is the first observational evidence of stellar-mass black holes weighing 25 solar masses or more.[154]
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+ Since then many more gravitational wave events have since been observed.[16]
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+ The proper motions of stars near the center of our own Milky Way provide strong observational evidence that these stars are orbiting a supermassive black hole.[155] Since 1995, astronomers have tracked the motions of 90 stars orbiting an invisible object coincident with the radio source Sagittarius A*. By fitting their motions to Keplerian orbits, the astronomers were able to infer, in 1998, that a 2.6 million M☉ object must be contained in a volume with a radius of 0.02 light-years to cause the motions of those stars.[156] Since then, one of the stars—called S2—has completed a full orbit. From the orbital data, astronomers were able to refine the calculations of the mass to 4.3 million M☉ and a radius of less than 0.002 light years for the object causing the orbital motion of those stars.[155] The upper limit on the object's size is still too large to test whether it is smaller than its Schwarzschild radius; nevertheless, these observations strongly suggest that the central object is a supermassive black hole as there are no other plausible scenarios for confining so much invisible mass into such a small volume.[156] Additionally, there is some observational evidence that this object might possess an event horizon, a feature unique to black holes.[157]
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+
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+ Due to conservation of angular momentum,[159] gas falling into the gravitational well created by a massive object will typically form a disk-like structure around the object. Artists' impressions such as the accompanying representation of a black hole with corona commonly depict the black hole as if it were a flat-space body hiding the part of the disk just behind it, but in reality gravitational lensing would greatly distort the image of the accretion disk.[160]
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+ Within such a disk, friction would cause angular momentum to be transported outward, allowing matter to fall further inward, thus releasing potential energy and increasing the temperature of the gas.[161]
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+ When the accreting object is a neutron star or a black hole, the gas in the inner accretion disk orbits at very high speeds because of its proximity to the compact object. The resulting friction is so significant that it heats the inner disk to temperatures at which it emits vast amounts of electromagnetic radiation (mainly X-rays). These bright X-ray sources may be detected by telescopes. This process of accretion is one of the most efficient energy-producing processes known; up to 40% of the rest mass of the accreted material can be emitted as radiation.[161] (In nuclear fusion only about 0.7% of the rest mass will be emitted as energy.) In many cases, accretion disks are accompanied by relativistic jets that are emitted along the poles, which carry away much of the energy. The mechanism for the creation of these jets is currently not well understood, in part due to insufficient data.[162]
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+ As such, many of the universe's more energetic phenomena have been attributed to the accretion of matter on black holes. In particular, active galactic nuclei and quasars are believed to be the accretion disks of supermassive black holes.[163] Similarly, X-ray binaries are generally accepted to be binary star systems in which one of the two stars is a compact object accreting matter from its companion.[163] It has also been suggested that some ultraluminous X-ray sources may be the accretion disks of intermediate-mass black holes.[164]
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+
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+ In November 2011 the first direct observation of a quasar accretion disk around a supermassive black hole was reported.[165][166]
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+
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+ X-ray binaries are binary star systems that emit a majority of their radiation in the X-ray part of the spectrum. These X-ray emissions are generally thought to result when one of the stars (compact object) accretes matter from another (regular) star. The presence of an ordinary star in such a system provides an opportunity for studying the central object and to determine if it might be a black hole.[163]
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+ If such a system emits signals that can be directly traced back to the compact object, it cannot be a black hole. The absence of such a signal does, however, not exclude the possibility that the compact object is a neutron star. By studying the companion star it is often possible to obtain the orbital parameters of the system and to obtain an estimate for the mass of the compact object. If this is much larger than the Tolman–Oppenheimer–Volkoff limit (the maximum mass a star can have without collapsing) then the object cannot be a neutron star and is generally expected to be a black hole.[163]
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+ The first strong candidate for a black hole, Cygnus X-1, was discovered in this way by Charles Thomas Bolton,[167] Louise Webster and Paul Murdin[168] in 1972.[169][170] Some doubt, however, remained due to the uncertainties that result from the companion star being much heavier than the candidate black hole. Currently, better candidates for black holes are found in a class of X-ray binaries called soft X-ray transients. In this class of system, the companion star is of relatively low mass allowing for more accurate estimates of the black hole mass. Moreover, these systems actively emit X-rays for only several months once every 10–50 years. During the period of low X-ray emission (called quiescence), the accretion disk is extremely faint allowing detailed observation of the companion star during this period. One of the best such candidates is V404 Cygni.[163]
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+ The X-ray emissions from accretion disks sometimes flicker at certain frequencies. These signals are called quasi-periodic oscillations and are thought to be caused by material moving along the inner edge of the accretion disk (the innermost stable circular orbit). As such their frequency is linked to the mass of the compact object. They can thus be used as an alternative way to determine the mass of candidate black holes.[171]
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+
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+ Astronomers use the term "active galaxy" to describe galaxies with unusual characteristics, such as unusual spectral line emission and very strong radio emission. Theoretical and observational studies have shown that the activity in these active galactic nuclei (AGN) may be explained by the presence of supermassive black holes, which can be millions of times more massive than stellar ones. The models of these AGN consist of a central black hole that may be millions or billions of times more massive than the Sun; a disk of gas and dust called an accretion disk; and two jets perpendicular to the accretion disk.[172][173]
163
+
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+ Although supermassive black holes are expected to be found in most AGN, only some galaxies' nuclei have been more carefully studied in attempts to both identify and measure the actual masses of the central supermassive black hole candidates. Some of the most notable galaxies with supermassive black hole candidates include the Andromeda Galaxy, M32, M87, NGC 3115, NGC 3377, NGC 4258, NGC 4889, NGC 1277, OJ 287, APM 08279+5255 and the Sombrero Galaxy.[175]
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+ It is now widely accepted that the center of nearly every galaxy, not just active ones, contains a supermassive black hole.[176] The close observational correlation between the mass of this hole and the velocity dispersion of the host galaxy's bulge, known as the M-sigma relation, strongly suggests a connection between the formation of the black hole and the galaxy itself.[177]
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+ Another way the black hole nature of an object may be tested in the future is through observation of effects caused by a strong gravitational field in their vicinity. One such effect is gravitational lensing: The deformation of spacetime around a massive object causes light rays to be deflected much as light passing through an optic lens. Observations have been made of weak gravitational lensing, in which light rays are deflected by only a few arcseconds. However, it has never been directly observed for a black hole.[179] One possibility for observing gravitational lensing by a black hole would be to observe stars in orbit around the black hole. There are several candidates for such an observation in orbit around Sagittarius A*.[179]
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+
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+ The evidence for stellar black holes strongly relies on the existence of an upper limit for the mass of a neutron star. The size of this limit heavily depends on the assumptions made about the properties of dense matter. New exotic phases of matter could push up this bound.[163] A phase of free quarks at high density might allow the existence of dense quark stars,[180] and some supersymmetric models predict the existence of Q stars.[181] Some extensions of the standard model posit the existence of preons as fundamental building blocks of quarks and leptons, which could hypothetically form preon stars.[182] These hypothetical models could potentially explain a number of observations of stellar black hole candidates. However, it can be shown from arguments in general relativity that any such object will have a maximum mass.[163]
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+ Since the average density of a black hole inside its Schwarzschild radius is inversely proportional to the square of its mass, supermassive black holes are much less dense than stellar black holes (the average density of a 108 M☉ black hole is comparable to that of water).[163] Consequently, the physics of matter forming a supermassive black hole is much better understood and the possible alternative explanations for supermassive black hole observations are much more mundane. For example, a supermassive black hole could be modelled by a large cluster of very dark objects. However, such alternatives are typically not stable enough to explain the supermassive black hole candidates.[163]
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+ The evidence for the existence of stellar and supermassive black holes implies that in order for black holes to not form, general relativity must fail as a theory of gravity, perhaps due to the onset of quantum mechanical corrections. A much anticipated feature of a theory of quantum gravity is that it will not feature singularities or event horizons and thus black holes would not be real artifacts.[183] For example, in the fuzzball model based on string theory, the individual states of a black hole solution do not generally have an event horizon or singularity, but for a classical/semi-classical observer the statistical average of such states appears just as an ordinary black hole as deduced from general relativity.[184]
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+ A few theoretical objects have been conjectured to match observations of astronomical black hole candidates identically or near-identically, but which function via a different mechanism. These include the gravastar, the black star,[185] and the dark-energy star.[186]
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+
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+ In 1971, Hawking showed under general conditions[Note 5] that the total area of the event horizons of any collection of classical black holes can never decrease, even if they collide and merge.[187] This result, now known as the second law of black hole mechanics, is remarkably similar to the second law of thermodynamics, which states that the total entropy of an isolated system can never decrease. As with classical objects at absolute zero temperature, it was assumed that black holes had zero entropy. If this were the case, the second law of thermodynamics would be violated by entropy-laden matter entering a black hole, resulting in a decrease of the total entropy of the universe. Therefore, Bekenstein proposed that a black hole should have an entropy, and that it should be proportional to its horizon area.[188]
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+ The link with the laws of thermodynamics was further strengthened by Hawking's discovery that quantum field theory predicts that a black hole radiates blackbody radiation at a constant temperature. This seemingly causes a violation of the second law of black hole mechanics, since the radiation will carry away energy from the black hole causing it to shrink. The radiation, however also carries away entropy, and it can be proven under general assumptions that the sum of the entropy of the matter surrounding a black hole and one quarter of the area of the horizon as measured in Planck units is in fact always increasing. This allows the formulation of the first law of black hole mechanics as an analogue of the first law of thermodynamics, with the mass acting as energy, the surface gravity as temperature and the area as entropy.[188]
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+ One puzzling feature is that the entropy of a black hole scales with its area rather than with its volume, since entropy is normally an extensive quantity that scales linearly with the volume of the system. This odd property led Gerard 't Hooft and Leonard Susskind to propose the holographic principle, which suggests that anything that happens in a volume of spacetime can be described by data on the boundary of that volume.[189]
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+
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+ Although general relativity can be used to perform a semi-classical calculation of black hole entropy, this situation is theoretically unsatisfying. In statistical mechanics, entropy is understood as counting the number of microscopic configurations of a system that have the same macroscopic qualities (such as mass, charge, pressure, etc.). Without a satisfactory theory of quantum gravity, one cannot perform such a computation for black holes. Some progress has been made in various approaches to quantum gravity. In 1995, Andrew Strominger and Cumrun Vafa showed that counting the microstates of a specific supersymmetric black hole in string theory reproduced the Bekenstein–Hawking entropy.[190] Since then, similar results have been reported for different black holes both in string theory and in other approaches to quantum gravity like loop quantum gravity.[191]
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+ Because a black hole has only a few internal parameters, most of the information about the matter that went into forming the black hole is lost. Regardless of the type of matter which goes into a black hole, it appears that only information concerning the total mass, charge, and angular momentum are conserved. As long as black holes were thought to persist forever this information loss is not that problematic, as the information can be thought of as existing inside the black hole, inaccessible from the outside, but represented on the event horizon in accordance with the holographic principle. However, black holes slowly evaporate by emitting Hawking radiation. This radiation does not appear to carry any additional information about the matter that formed the black hole, meaning that this information appears to be gone forever.[192]
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+
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+ The question whether information is truly lost in black holes (the black hole information paradox) has divided the theoretical physics community (see Thorne–Hawking–Preskill bet). In quantum mechanics, loss of information corresponds to the violation of a property called unitarity, and it has been argued that loss of unitarity would also imply violation of conservation of energy,[193] though this has also been disputed.[194] Over recent years evidence has been building that indeed information and unitarity are preserved in a full quantum gravitational treatment of the problem.[195]
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+ One attempt to resolve the black hole information paradox is known as black hole complementarity. In 2012, the "firewall paradox" was introduced with the goal of demonstrating that black hole complementarity fails to solve the information paradox. According to quantum field theory in curved spacetime, a single emission of Hawking radiation involves two mutually entangled particles. The outgoing particle escapes and is emitted as a quantum of Hawking radiation; the infalling particle is swallowed by the black hole. Assume a black hole formed a finite time in the past and will fully evaporate away in some finite time in the future. Then, it will emit only a finite amount of information encoded within its Hawking radiation. According to research by physicists like Don Page[196][197] and Leonard Susskind, there will eventually be a time by which an outgoing particle must be entangled with all the Hawking radiation the black hole has previously emitted. This seemingly creates a paradox: a principle called "monogamy of entanglement" requires that, like any quantum system, the outgoing particle cannot be fully entangled with two other systems at the same time; yet here the outgoing particle appears to be entangled both with the infalling particle and, independently, with past Hawking radiation.[198] In order to resolve this contradiction, physicists may eventually be forced to give up one of three time-tested principles: Einstein's equivalence principle, unitarity, or local quantum field theory. One possible solution, which violates the equivalence principle, is that a "firewall" destroys incoming particles at the event horizon.[199] In general, which if any of these assumptions should be abandoned remains a topic of debate.[194]
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+ A pig is any of the animals in the genus Sus, within the even-toed ungulate family Suidae. Pigs include domestic pigs and their ancestor, the common Eurasian wild boar (Sus scrofa), along with other species. Pigs, like all suids, are native to the Eurasian and African continents, ranging from Europe to the Pacific islands. Suids other than the pig are the babirusa of Indonesia, the pygmy hog of Asia, the warthog of Africa, and another genus of pigs from Africa. The suids are a sister clade to peccaries.
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+ Juvenile pigs are known as piglets.[1] Pigs are highly social and intelligent animals.[2]
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+ With around 1 billion individuals alive at any time, the domestic pig is among the most populous large mammals in the world.[3][4] Pigs are omnivores and can consume a wide range of food.[5] Pigs are biologically similar to humans and are thus frequently used for human medical research.[6]
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+ The Online Etymology Dictionary provides anecdotal evidence as well as linguistic, saying that the term derives
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+
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+ probably from Old English *picg, found in compounds, ultimate origin unknown. Originally "young pig" (the word for adults was swine). Apparently related to Low German bigge, Dutch big ("but the phonology is difficult" -- OED). ... Another Old English word for "pig" was fearh, related to furh "furrow," from PIE *perk- "dig, furrow" (source also of Latin porc-us "pig," see pork). "This reflects a widespread IE tendency to name animals from typical attributes or activities" [Roger Lass]. Synonyms grunter, oinker are from sailors' and fishermen's euphemistic avoidance of uttering the word pig at sea, a superstition perhaps based on the fate of the Gadarene swine, who drowned.[7]
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+
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+ The Online Etymology Dictionary also traces the evolution of sow, the term for a female pig, through various historical languages:
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+ Old English sugu, su "female of the swine," from Proto-Germanic *su- (cognates: Old Saxon, Old High German su, German Sau, Dutch zeug, Old Norse syr), from PIE root *su- (cognates: Sanskrit sukarah "wild boar, swine;" Avestan hu "wild boar;" Greek hys "swine;" Latin sus "swine", suinus "pertaining to swine"; Old Church Slavonic svinija "swine;" Lettish sivens "young pig;" Welsh hucc, Irish suig "swine; Old Irish socc "snout, plowshare"), possibly imitative of pig noise; note that Sanskrit sukharah means "maker of (the sound) su.[7]
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+ An adjectival form is porcine. Another adjectival form (technically for the subfamily rather than genus name) is suine (comparable to bovine, canine, etc.); for the family, it is suid (as with bovid, canid).
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+ A typical pig has a large head with a long snout that is strengthened by a special prenasal bone and by a disk of cartilage at the tip.[8] The snout is used to dig into the soil to find food and is a very acute sense organ. There are four hoofed toes on each foot, with the two larger central toes bearing most of the weight, but the outer two also being used in soft ground.[9]
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+ The dental formula of adult pigs is 3.1.4.33.1.4.3, giving a total of 44 teeth. The rear teeth are adapted for crushing. In the male, the canine teeth form tusks, which grow continuously and are sharpened by constantly being ground against each other.[8]
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+ Occasionally, captive mother pigs may savage their own piglets, often if they become severely stressed.[10] Some attacks on newborn piglets are non-fatal. Others may cause the death of the piglets and sometimes, the mother may eat the piglets. It is estimated that 50% of piglet fatalities are due to the mother attacking, or unintentionally crushing, the newborn pre-weaned animals.[11]
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+ With around 1 billion individuals alive at any time, the domestic pig is one of the most numerous large mammals on the planet.[3][4]
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+ The ancestor of the domestic pig is the wild boar, which is one of the most numerous and widespread large mammals. Its many subspecies are native to all but the harshest climates of continental Eurasia and its islands and Africa as well, from Ireland and India to Japan and north to Siberia.
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+ Long isolated from other pigs on the many islands of Indonesia, Malaysia, and the Philippines, pigs have evolved into many different species, including wild boar, bearded pigs, and warty pigs. Humans have introduced pigs into Australia, North and South America, and numerous islands, either accidentally as escaped domestic pigs which have gone feral, or as wild boar.
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+ The wild boar (Sus scrofa) can take advantage of any forage resources. Therefore, they can live in virtually any productive habitat that can provide enough water to sustain large mammals such as pigs. If there is increased foraging of wild boars in certain areas, they can cause a nutritional shortage which can cause the pig population to decrease. If the nutritional state returns to normal, the pig population will most likely rise due to the pigs' naturally increased reproduction rate.[12]
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+ Pigs are omnivores, which means that they consume both plants and animals. In the wild, they are foraging animals, primarily eating leaves, roots, fruits, and flowers, in addition to some insects and fish. As livestock, pigs are fed mostly corn and soybean meal[13] with a mixture of vitamins and minerals added to the diet. Traditionally, they were raised on dairy farms and called "mortgage lifters", due to their ability to use the excess milk as well as whey from cheese and butter making combined with pasture.[14] Older pigs will consume three to five gallons of water per day.[15] When kept as pets, the optimal healthy diet consists mainly of a balanced diet of raw vegetables, although some may give their pigs conventional mini pig pellet feed.[16]
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+ Domesticated pigs, especially miniature breeds, are commonly kept as pets.[17] Domestic pigs are raised commercially as livestock; materials that are garnered include their meat (known as pork), leather, and their bristly hairs which are used to make brushes. Because of their foraging abilities and excellent sense of smell, they are used to find truffles in many European countries. Both wild and feral pigs are commonly hunted.
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+ The relatively short, stiff, coarse hairs of the pig are called bristles, and were once so commonly used in paintbrushes that in 1946 the Australian Government launched Operation Pig Bristle. In May 1946, in response to a shortage of pig bristles for paintbrushes to paint houses in the post-World War II construction boom, the Royal Australian Air Force (RAAF) flew in 28 short tons of pig bristles from China, their only commercially available source at the time.[18]
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+ Human skin is very similar to pig skin, therefore pig skin has been used in many preclinical studies.[19][20] In addition to providing use in biomedical research[19][20] and for drug testing,[21] genetic advances in human healthcare have provided a pathway for domestic pigs to become xenotransplantation candidates
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+ for humans.[22]
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+ The genus Sus is currently thought to contain eight living species. A number of extinct species (†) are known from fossils.
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+ The pygmy hog, formerly Sus salvanius is now placed in the monotypic genus Porcula.[23]
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+ Pigs have been domesticated since ancient times in the Old World. Archaeological evidence suggests that pigs were being managed in the wild in a way similar to the way they are managed by some modern New Guineans from wild boar as early as 13,000–12,700 BP in the Near East in the Tigris Basin,[24] Çayönü, Cafer Höyük, Nevalı Çori.[25] Remains of pigs have been dated to earlier than 11,400 BP in Cyprus that must have been introduced from the mainland which suggests domestication in the adjacent mainland by then.[26] A separate domestication also occurred in China.[27]
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+ In India, pigs have been domesticated for a long time mostly in Goa and some rural areas for pig toilets. This was also done in China. Though ecologically logical as well as economical, pig toilets are waning in popularity as use of septic tanks and/or sewerage systems is increasing in rural areas.
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+ Pigs were brought to southeastern North America from Europe by Hernando de Soto and other early Spanish explorers. Pigs are particularly valued in China and on certain oceanic islands, where their self-sufficiency allows them to be turned loose, although the practice is not without its drawbacks (see environmental impact).
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+ The domestic pig (Sus scrofa domesticus) is usually given the scientific name Sus scrofa, although some taxonomists call it S. domesticus, reserving S. scrofa for the wild boar. It was domesticated approximately 5,000 to 7,000 years ago. The upper canines form sharp distinctive tusks that curve outward and upward. Compared to other artiodactyles, their head is relatively long, pointed, and free of warts. Their head and body length ranges from 0.9 to 1.8 m (35 to 71 in) and they can weigh between 50 and 350 kg (110 and 770 lb).
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+ In November 2012, scientists managed to sequence the genome of the domestic pig. The similarities between the pig and human genomes mean that the new data may have wide applications in the study and treatment of human genetic diseases.[28][29][30]
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+ In August 2015, a study looked at over 100 pig genome sequences to ascertain their process of domestication. The process of domestication was assumed to have been initiated by humans, involved few individuals and relied on reproductive isolation between wild and domestic forms. The study found that the assumption of reproductive isolation with population bottlenecks was not supported. The study indicated that pigs were domesticated separately in Western Asia and China, with Western Asian pigs introduced into Europe where they crossed with wild boar. A model that fitted the data included admixture with a now extinct ghost population of wild pigs during the Pleistocene. The study also found that despite back-crossing with wild pigs, the genomes of domestic pigs have strong signatures of selection at DNA loci that affect behavior and morphology. The study concluded that human selection for domestic traits likely counteracted the homogenizing effect of gene flow from wild boars and created domestication islands in the genome. The same process may also apply to other domesticated animals.[31]
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+ [32]
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+ Pigs have been important in culture across the world since neolithic times. They appear in art, literature, and religion. In Asia the wild boar is one of 12 animal images comprising the Chinese zodiac, while in Europe the boar represents a standard charge in heraldry. In Islam and Judaism pigs and those who handle them are viewed negatively, and the consumption of pork is forbidden.[33][34] Pigs are alluded to in animal epithets and proverbs.[35][36]
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+ The pig has been celebrated throughout Europe since ancient times in its carnivals, the name coming from the Italian carne levare, the lifting of meat.[37]
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+ Pigs have been brought into literature for varying reasons, ranging from the pleasures of eating, as in Charles Lamb's A Dissertation upon Roast Pig, to William Golding's Lord of the Flies (with the fat character "Piggy"), where the rotting boar's head on a stick represents Beelzebub, "lord of the flies" being the direct translation of the Hebrew בעל זבוב, and George Orwell's allegorical novel Animal Farm, where the central characters, representing Soviet leaders, are all pigs.[38][39][40][37]
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+ Domestic pigs that have escaped from urban areas or were allowed to forage in the wild, and in some cases wild boars which were introduced as prey for hunting, have given rise to large populations of feral pigs in North and South America, Australia, New Zealand, Hawaii, and other areas where pigs are not native. Accidental or deliberate releases of pigs into countries or environments where they are an alien species have caused extensive environmental change. Their omnivorous diet, aggressive behaviour, and their feeding method of rooting in the ground all combine to severely alter ecosystems unused to pigs. Pigs will even eat small animals and destroy nests of ground nesting birds.[8] The Invasive Species Specialist Group lists feral pigs on the list of the world's 100 worst invasive species and says:[41]
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+ Feral pigs like other introduced mammals are major drivers of extinction and ecosystem change. They have been introduced into many parts of the world, and will damage crops and home gardens as well as potentially spreading disease. They uproot large areas of land, eliminating native vegetation and spreading weeds. This results in habitat alteration, a change in plant succession and composition and a decrease in native fauna dependent on the original habitat.
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+ Because of the biological similarities between each other, pigs can harbour a range of parasites and diseases that can be transmitted to humans. These include trichinosis, Taenia solium, cysticercosis, and brucellosis. Pigs are also known to host large concentrations of parasitic ascarid worms in their digestive tract.[42]
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+ Some strains of influenza are endemic in pigs. Pigs also can acquire human influenza.[further explanation needed]