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+In religion and folklore, Hell is an afterlife location in which evil souls are subjected to punitive suffering, often torture as eternal punishment after death. Religions with a linear divine history often depict hells as eternal destinations, the biggest examples of which are Christianity and Islam, whereas religions with reincarnation usually depict a hell as an intermediary period between incarnations, as is the case in the dharmic religions. Religions typically locate hell in another dimension or under Earth's surface. Other afterlife destinations include Heaven, Paradise, Purgatory, Limbo, and the underworld.
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+Other religions, which do not conceive of the afterlife as a place of punishment or reward, merely describe an abode of the dead, the grave, a neutral place that is located under the surface of Earth (for example, see Kur, Hades, and Sheol). Such places are sometimes equated with the English word hell, though a more correct translation would be "underworld" or "world of the dead". The ancient Mesopotamian, Greek, Roman, and Finnic religions include entrances to the underworld from the land of the living.
+
+The modern English word hell is derived from Old English hel, helle (first attested around 725 AD to refer to a nether world of the dead) reaching into the Anglo-Saxon pagan period.[1] The word has cognates in all branches of the Germanic languages, including Old Norse hel (which refers to both a location and goddess-like being in Norse mythology), Old Frisian helle, Old Saxon hellia, Old High German hella, and Gothic halja. All forms ultimately derive from the reconstructed Proto-Germanic feminine noun *xaljō or *haljō ('concealed place, the underworld'). In turn, the Proto-Germanic form derives from the o-grade form of the Proto-Indo-European root *kel-, *kol-: 'to cover, conceal, save'.[2] Indo-European cognates include Latin cēlāre ("to hide", related to the English word cellar) and early Irish ceilid ("hides"). Upon the Christianization of the Germanic peoples, extension of Proto-Germanic *xaljō were reinterpreted to denote the underworld in Christian mythology,[1][3] for which see Gehenna.
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+Related early Germanic terms and concepts include Proto-Germanic *xalja-rūnō(n), a feminine compound noun, and *xalja-wītjan, a neutral compound noun. This form is reconstructed from the Latinized Gothic plural noun *haliurunnae (attested by Jordanes; according to philologist Vladimir Orel, meaning 'witches'), Old English helle-rúne ('sorceress, necromancer', according to Orel), and Old High German helli-rūna 'magic'. The compound is composed of two elements: *xaljō (*haljō) and *rūnō, the Proto-Germanic precursor to Modern English rune.[4] The second element in the Gothic haliurunnae may however instead be an agent noun from the verb rinnan ("to run, go"), which would make its literal meaning "one who travels to the netherworld".[5][6]
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+Proto-Germanic *xalja-wītjan (or *halja-wītjan) is reconstructed from Old Norse hel-víti 'hell', Old English helle-wíte 'hell-torment, hell', Old Saxon helli-wīti 'hell', and the Middle High German feminine noun helle-wīze. The compound is a compound of *xaljō (discussed above) and *wītjan (reconstructed from forms such as Old English witt 'right mind, wits', Old Saxon gewit 'understanding', and Gothic un-witi 'foolishness, understanding').[7]
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+Hell appears in several mythologies and religions. It is commonly inhabited by demons and the souls of dead people. A fable about Hell which recurs in folklore across several cultures is the allegory of the long spoons. Hell is often depicted in art and literature, perhaps most famously in Dante's Divine Comedy.
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+Punishment in Hell typically corresponds to sins committed during life. Sometimes these distinctions are specific, with damned souls suffering for each sin committed (see for example Plato's myth of Er or Dante's The Divine Comedy), but sometimes they are general, with condemned sinners relegated to one or more chamber of Hell or to a level of suffering.
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+In many religious cultures, including Christianity and Islam, Hell is often depicted as fiery, painful, and harsh, inflicting suffering on the guilty. Despite these common depictions of Hell as a place of fire, some other traditions portray Hell as cold. Buddhist – and particularly Tibetan Buddhist – descriptions of Hell feature an equal number of hot and cold Hells. Among Christian descriptions Dante's Inferno portrays the innermost (9th) circle of Hell as a frozen lake of blood and guilt.[11]
+But cold also played a part in earlier Christian depictions of Hell, beginning with the Apocalypse of Paul, originally from the early third century;[12] the "Vision of Dryhthelm" by the Venerable Bede from the seventh century;[13] "St Patrick's Purgatory", "The Vision of Tundale" or "Visio Tnugdali", and the "Vision of the Monk of Eynsham", all from the twelfth century;[14] and the "Vision of Thurkill" from the early thirteenth century.[15]
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+The Sumerian afterlife was a dark, dreary cavern located deep below the ground,[16] where inhabitants were believed to continue "a shadowy version of life on earth".[16] This bleak domain was known as Kur,[17]:114 and was believed to be ruled by the goddess Ereshkigal.[16][18]:184 All souls went to the same afterlife,[16] and a person's actions during life had no effect on how the person would be treated in the world to come.[16]
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+The souls in Kur were believed to eat nothing but dry dust[17]:58 and family members of the deceased would ritually pour libations into the dead person's grave through a clay pipe, thereby allowing the dead to drink.[17]:58 Nonetheless, funerary evidence indicates that some people believed that the goddess Inanna, Ereshkigal's younger sister, had the power to award her devotees with special favors in the afterlife.[16][19] During the Third Dynasty of Ur, it was believed that a person's treatment in the afterlife depended on how he or she was buried;[17]:58 those that had been given sumptuous burials would be treated well,[17]:58 but those who had been given poor burials would fare poorly.[17]:58
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+The entrance to Kur was believed to be located in the Zagros mountains in the far east.[17]:114 It had seven gates, through which a soul needed to pass.[16] The god Neti was the gatekeeper.[18]:184[17]:86 Ereshkigal's sukkal, or messenger, was the god Namtar.[17]:134[18]:184 Galla were a class of demons that were believed to reside in the underworld;[17]:85 their primary purpose appears to have been to drag unfortunate mortals back to Kur.[17]:85 They are frequently referenced in magical texts,[17]:85–86 and some texts describe them as being seven in number.[17]:85–86 Several extant poems describe the galla dragging the god Dumuzid into the underworld.[17]:86 The later Mesopotamians knew this underworld by its East Semitic name: Irkalla. During the Akkadian Period, Ereshkigal's role as the ruler of the underworld was assigned to Nergal, the god of death.[16][18]:184 The Akkadians attempted to harmonize this dual rulership of the underworld by making Nergal Ereshkigal's husband.[16]
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+With the rise of the cult of Osiris during the Middle Kingdom the "democratization of religion" offered to even his humblest followers the prospect of eternal life, with moral fitness becoming the dominant factor in determining a person's suitability. At death a person faced judgment by a tribunal of forty-two divine judges. If they had led a life in conformance with the precepts of the goddess Maat, who represented truth and right living, the person was welcomed into the heavenly reed fields. If found guilty the person was thrown to Ammit, the "devourer of the dead" and would be condemned to the lake of fire.[21] The person taken by the devourer is subject first to terrifying punishment and then annihilated. These depictions of punishment may have influenced medieval perceptions of the inferno in hell via early Christian and Coptic texts.[22] Purification for those considered justified appears in the descriptions of "Flame Island", where humans experience the triumph over evil and rebirth. For the damned complete destruction into a state of non-being awaits but there is no suggestion of eternal torture; the weighing of the heart in Egyptian mythology can lead to annihilation.[23][24] The Tale of Khaemwese describes the torment of a rich man, who lacked charity, when he dies and compares it to the blessed state of a poor man who has also died.[25]
+Divine pardon at judgement always remained a central concern for the ancient Egyptians.[26]
+
+Modern understanding of Egyptian notions of hell relies on six ancient texts:[27]
+
+In classic Greek mythology, below Heaven, Earth, and Pontus is Tartarus, or Tartaros (Greek Τάρταρος, deep place). It is either a deep, gloomy place, a pit or abyss used as a dungeon of torment and suffering that resides within Hades (the entire underworld) with Tartarus being the hellish component. In the Gorgias, Plato (c. 400 BC) wrote that souls of the deceased were judged after they payed for crossing the river of the dead and those who received punishment were sent to Tartarus.[28] As a place of punishment, it can be considered a hell. The classic Hades, on the other hand, is more similar to Old Testament Sheol. The Romans later adopted these views.
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+The hells of Europe include Breton mythology's "Anaon", Celtic mythology's "Uffern", Slavic mythology's "Peklo", the hell of Sami mythology and Finnish "tuonela" ("manala").
+
+The hells of Asia include the Bagobo "Gimokodan" (which is believed to be more of an otherworld, where the Red Region is reserved who those who died in battle, while ordinary people go to the White Region)[29] and ancient Indian mythology's "Kalichi" or "Naraka".
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+According to a few sources, hell is below ground, and described as an uninviting wet[30] or fiery place reserved for sinful people in the Ainu religion, as stated by missionary John Batchelor.[31] However, belief in hell does not appear in oral tradition of the Ainu.[32] Instead, there is belief within the Ainu religion that the soul of the deceased (ramat) would become a kamuy after death.[32] There is also belief that the soul of someone who has been wicked during lifetime, committed suicide, got murdered or died in great agony would become a ghost (tukap) who would haunt the living,[32] to come to fulfillment from which it was excluded during life.[33]
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+In Tengrism, it was believed that the wicked would get punished in Tamag before they would be brought to the third floor of the sky.[34]
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+In Taoism, hell is represented by Diyu.
+
+The Hell of Swahili mythology is called kuzimu, and believe in it deleveloped in the 7th and 8th century under the influence of Muslim merchants at the east African coast.[35] It is imagined as a very cold place.[35] Serer religion rejects the general notion of heaven and hell.[36] In Serer religion, acceptance by the ancestors who have long departed is as close to any heaven as one can get. Rejection and becoming a wandering soul is a sort of hell for one passing over. The souls of the dead must make their way to Jaaniw (the sacred dwelling place of the soul). Only those who have lived their lives on earth in accordance with Serer doctrines will be able to make this necessary journey and thus accepted by the ancestors. Those who can't make the journey become lost and wandering souls, but they do not burn in "hell fire".[36][37]
+
+The hells of the Americas include the Aztec religion's Mictlan, Inuit religion's Adlivun, and the Yanomami religion's Shobari Waka. In Mayan religion, Xibalba (or Metnal) is the dangerous underworld of nine levels. The road into and out of it is said to be steep, thorny and very forbidding. Ritual healers would intone healing prayers banishing diseases to Xibalba. Much of the Popol Vuh describes the adventures of the Maya Hero Twins in their cunning struggle with the evil lords of Xibalba.
+
+The Aztecs believed that the dead traveled to Mictlan, a neutral place found far to the north. There was also a legend of a place of white flowers, which was always dark, and was home to the gods of death, particularly Mictlantecutli and his spouse Mictlantecihuatl, which means literally "lords of Mictlan". The journey to Mictlan took four years, and the travelers had to overcome difficult tests, such as passing a mountain range where the mountains crashed into each other, a field where the wind carried flesh-scraping knives, and a river of blood with fearsome jaguars.
+
+In pre-Christian Fijian mythology there was belief in an underworld called Murimuria.
+
+Hell is conceived of in most Abrahamic religions as a place of, or a form of, punishment.[38]
+
+Early Judaism had no concept of Hell, although the concept of an afterlife was introduced during the Hellenistic period, apparently from neighboring Hellenistic religions. It occurs for example in the Book of Daniel. Daniel 12:2 proclaims "And many of those who sleep in the dust of the earth shall awake, Some to everlasting life, Some to shame and everlasting contempt."
+
+Judaism does not have a specific doctrine about the afterlife, but it does have a mystical/Orthodox tradition of describing Gehinnom. Gehinnom is not Hell, but originally a grave and in later times a sort of Purgatory where one is judged based on one's life's deeds, or rather, where one becomes fully aware of one's own shortcomings and negative actions during one's life. The Kabbalah explains it as a "waiting room" (commonly translated as an "entry way") for all souls (not just the wicked). The overwhelming majority of rabbinic thought maintains that people are not in Gehinnom forever; the longest that one can be there is said to be 12 months, however there has been the occasional noted exception. Some consider it a spiritual forge where the soul is purified for its eventual ascent to Olam Habah (heb. עולם הבא; lit. "The world to come", often viewed as analogous to heaven). This is also mentioned in the Kabbalah, where the soul is described as breaking, like the flame of a candle lighting another: the part of the soul that ascends being pure and the "unfinished" piece being reborn.
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+According to Jewish teachings, hell is not entirely physical; rather, it can be compared to a very intense feeling of shame. People are ashamed of their misdeeds and this constitutes suffering which makes up for the bad deeds. When one has so deviated from the will of God, one is said to be in Gehinnom. This is not meant to refer to some point in the future, but to the very present moment. The gates of teshuva (return) are said to be always open, and so one can align his will with that of God at any moment. Being out of alignment with God's will is itself a punishment according to the Torah.
+
+Many scholars of Jewish mysticism, particularly of the Kabbalah, describe seven "compartments" or "habitations" of Hell, just as they describe seven divisions of Heaven. These divisions go by many different names, and the most frequently mentioned are as follows:[39]
+
+Besides those mentioned above, there also exist additional terms that have been often used to either refer to Hell in general or to some region of the underworld:
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+For more information, see Qliphoth.
+
+Maimonides declares in his 13 principles of faith that the hells of the rabbinic literature were pedagocically motivated inventions to encourage respect of the Torah commandements by mankind, which had been regarded as immature.[49] Instead of being sent to hell, the souls of the wicked would actually get annihilated.[50]
+
+The Christian doctrine of hell derives from passages in the New Testament. The word hell does not appear in the Greek New Testament; instead one of three words is used: the Greek words Tartarus or Hades, or the Hebrew word Gehinnom.
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+In the Septuagint and New Testament the authors used the Greek term Hades for the Hebrew Sheol, but often with Jewish rather than Greek concepts in mind. In the Jewish concept of Sheol, such as expressed in Ecclesiastes,[51] Sheol or Hades is a place where there is no activity. However, since Augustine, some[which?] Christians have believed that the souls of those who die either rest peacefully, in the case of Christians, or are afflicted, in the case of the damned, after death until the resurrection.[52]
+
+While these three terms are translated in the KJV as "hell" these three terms have three very different meanings.
+
+The Roman Catholic Church defines Hell as "a state of definitive self-exclusion from communion with God and the blessed." One finds oneself in Hell as the result of dying in mortal sin without repenting and accepting God's merciful love, becoming eternally separated from him by one's own free choice[66] immediately after death.[67] In the Roman Catholic Church, many other Christian churches, such as the Baptists and Episcopalians, and some Greek Orthodox churches,[68] Hell is taught as the final destiny of those who have not been found worthy after the general resurrection and last judgment,[69][70][71] where they will be eternally punished for sin and permanently separated from God. The nature of this judgment is inconsistent with many Protestant churches teaching the saving comes from accepting Jesus Christ as their savior, while the Greek Orthodox and Catholic Churches teach that the judgment hinges on both faith and works. However, many Liberal Christians throughout Liberal Protestant and Anglican churches believe in universal reconciliation (see below), even though it contradicts the traditional doctrines that are usually held by the evangelicals within their denominations.[72]
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+Some modern Christian theologians subscribe to the doctrines of conditional immortality. Conditional immortality is the belief that the soul dies with the body and does not live again until the resurrection. As with other Jewish writings of the Second Temple period, the New Testament text distinguishes two words, both translated "Hell" in older English Bibles: Hades, "the grave", and Gehenna where God "can destroy both body and soul".[73] A minority of Christians read this to mean that neither Hades nor Gehenna are eternal but refer to the ultimate destruction of the wicked in the Lake of Fire in a consuming fire after resurrection. However, because of the Greek words used in translating from the Hebrew text, the Hebrew ideas have become confused with Greek myths and ideas.  In the Hebrew text when people died they went to Sheol, the grave[74] and the wicked ultimately went to Gehenna and were consumed by fire. The Hebrew words for "the grave" or "death" or "eventual destruction of the wicked", were translated using Greek words and later texts became a mix of mistranslation, pagan influence, and Greek myth.[75]
+
+Christian mortalism is the doctrine that all men and women, including Christians, must die, and do not continue and are not conscious after death. Therefore, annihilationism includes the doctrine that "the wicked" are also destroyed rather than tormented forever in traditional "Hell" or the lake of fire. Christian mortalism and annihilationism are directly related to the doctrine of conditional immortality, the idea that a human soul is not immortal unless it is given eternal life at the second coming of Christ and resurrection of the dead.
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+Biblical scholars looking at the issue through the Hebrew text have denied the teaching of innate immortality.[76][77] Rejection of the immortality of the soul, and advocacy of Christian mortalism, was a feature of Protestantism since the early days of the Reformation with Martin Luther himself rejecting the traditional idea, though his mortalism did not carry into orthodox Lutheranism. One of the most notable English opponents of the immortality of the soul was Thomas Hobbes who describes the idea as a Greek "contagion" in Christian doctrine.[78] Modern proponents of conditional immortality include some in the Anglican church such as N.T. Wright[79] and as denominations the Seventh-day Adventists, Bible Students, Jehovah's Witnesses, Christadelphians, Living Church of God, The Church of God International, and some other Protestant Christians, as well as recent Roman Catholic teaching. It is not Roman Catholic dogma that anyone is in Hell,[80] though many individual Catholics do not share this view. The 1993 Catechism of the Catholic Church states:[81] "This state of definitive self-exclusion from communion with God and the blessed is called 'hell'" and[82] "... they suffer the punishments of hell, "eternal fire." The chief punishment of hell is eternal separation from God" (CCC 1035).  During an Audience in 1999, Pope John Paul II commented: "images of hell that Sacred Scripture presents to us must be correctly interpreted. They show the complete frustration and emptiness of life without God. Rather than a place, hell indicates the state of those who freely and definitively separate themselves from God, the source of all life and joy."[83]
+
+The Seventh-day Adventist Church's official beliefs support annihilationism.[84][85] They deny the Catholic purgatory and teach that the dead lie in the grave until they are raised for a last judgment, both the righteous and wicked await the resurrection at the Second Coming. Seventh-day Adventists believe that death is a state of unconscious sleep until the resurrection. They base this belief on biblical texts such as Ecclesiastes 9:5 which states "the dead know nothing", and 1 Thessalonians 4:13–18 which contains a description of the dead being raised from the grave at the second coming. These verses, it is argued, indicate that death is only a period or form of slumber.
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+Adventists teach that the resurrection of the righteous will take place shortly after the second coming of Jesus, as described in Revelation 20:4–6 that follows Revelation 19:11–16, whereas the resurrection of the wicked will occur after the millennium, as described in Revelation 20:5 and 20:12–13 that follow Revelation 20:4 and 6–7, though Revelation 20:12–13 and 15 actually describe a mixture of saved and condemned people being raised from the dead and judged. Adventists reject the traditional doctrine of hell as a state of everlasting conscious torment, believing instead that the wicked will be permanently destroyed after the millennium by the lake of fire, which is called 'the second death' in Revelation 20:14.
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+Those Adventist doctrines about death and hell reflect an underlying belief in: (a) conditional immortality (or conditionalism), as opposed to the immortality of the soul; and (b) the monistic nature of human beings, in which the soul is not separable from the body, as opposed to bipartite or tripartite conceptions, in which the soul is separable.
+
+Jehovah's Witnesses hold that the soul ceases to exist when the person dies[86] and therefore that Hell (Sheol or Hades) is a state of non-existence.[86] In their theology, Gehenna differs from Sheol or Hades in that it holds no hope of a resurrection.[86] Tartarus is held to be the metaphorical state of debasement of the fallen angels between the time of their moral fall (Genesis chapter 6) until their post-millennial destruction along with Satan (Revelation chapter 20).[87]
+
+Bible Students and Christadelphians also believe in annihilationism.
+
+Christian Universalists believe in universal reconciliation, the belief that all human souls will be eventually reconciled with God and admitted to Heaven.[88] This belief is held by some Unitarian-Universalists.[89][90][91]
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+According to Emanuel Swedenborg's Second Coming Christian revelation, hell exists because evil people want it.[92] They, not God, introduced evil to the human race.[93] In Swedenborgianism, every soul joins the like-minded group after death in which it feels the most comfortable. Hell is therefore believed to be a place of hapiness for the souls which delight in evilness.[94]
+
+Members of The Church of Jesus Christ of Latter-day Saints (LDS Church) teach that hell is a state between death and resurrection, in which those spirits who did not repent while on earth must suffer for their own sins (Doctrine and Covenants 19:15–17[95]). After that, only the Sons of perdition, who committed the Eternal sin, would be cast into Outer darkness. However, according to Mormon faith, committing the Eternal sin requires so much knowledge that most persons cannot do this.[96] Satan and Cain are counted as examples of Sons of perdition.
+
+In Islam, jahannam (in Arabic: جهنم) (related to the Hebrew word gehinnom) is the counterpart to heaven and likewise divided into seven layers, both co-existing with the temporal world,[97] filled with blazing fire, boiling water, and a variety of other torments for those who have been condemned to it in the hereafter. In the Quran, God declares that the fire of Jahannam is prepared for both mankind and jinn.[98][99] After the Day of Judgement, it is to be occupied by those who do not believe in God, those who have disobeyed his laws, or rejected his messengers.[100] "Enemies of Islam" are sent to Hell immediately upon their deaths.[101] Muslim modernists downplay the vivid descriptions of hell common during Classical period, on one hand reaffirming that the afterlife must not be denied, but simultaneously asserting its exact nature remains unknown. Other modern Muslims continue the line of Sufism as an interiorized hell, combining the eschatological thoughts of Ibn Arabi and Rumi with Western philosophy.[102] Although disputed by some scholars, most scholars consider jahannam to be eternal.[103][104] There is belief that the fire which represents the own bad deeds can already be seen during the Punishment of the Grave, and that the spiritual pain caused by this can lead to purification of the soul.[105]
+
+Medieval sources usually identified hell with the seven layers of the earth mentioned in Surah 65:12, inhabited by devils, harsh angels, scorpions and serpents, who torment the sinners. They described thorny shrubs, seas filled with blood and fire and darkness only illuminated by the flames of hell.[106] However, some sources also mention a place of extreme cold at the bottom of hell, called Zamhareer, characterized in as being unbearably cold, with blizzards, ice, and snow.[107] Maalik is thought of as the keeper of the gates of hell, namely appears in Ibn Abbas' Isra and Mi'raj.[108] Over hell, a narrow bridge called As-Sirāt is spanned. On Judgement Day one must pass over it to reach paradise, but those destined for hell will find too narrow and fall from into their new abode.[109] Iblis, the temporary ruler of hell,[110] is thought of residing in the bottom of hell, from where he commands his hosts of infernal demons.[111][112] But contrary to Christian traditions, Iblis and his infernal hosts do not wage war against God,[113] his enmity applies against humanity only. Further, his dominion in hell is also his punishment. According to the Muwatta Hadith, the Bukhari Hadith, the Tirmidhi Hadith, and the Kabir Hadith, Muhammad claimed that the fire of Jahannam is not red, but pitch-black, and is 70 times hotter than ordinary fire, and is much more painful than ordinary fire.[citation needed]
+
+Polytheism (shirk) is regarded as a particularly grievous sin; therefore entering Paradise is forbidden to a polytheist (musyrik) because his place is Hell;[114] and the lowest pit of Hell (Hawiyah), is intended for hypocrites who claimed aloud to believe in God and his messenger but in their hearts did not.[115] Not all Muslims and scholars agree whether hell is an eternal destination or whether some or all of the condemned will eventually be forgiven and allowed to enter paradise.[101][113][116][117]
+
+In the Bahá'í Faith, the conventional descriptions of Hell and Heaven are considered to be symbolic representations of spiritual conditions. The Bahá'í writings describe closeness to God to be Heaven, and conversely, remoteness from God as Hell.[118] The Bahá'í writings state that the soul is immortal and after death it will continue to progress until it finally attains God's presence.[119]
+
+In "Devaduta Sutta", the 130th discourse of the Majjhima Nikaya, Buddha teaches about hell in vivid detail. Buddhism teaches that there are five[citation needed] (sometimes six[citation needed]) realms of rebirth, which can then be further subdivided into degrees of agony or pleasure. Of these realms, the hell realms, or Naraka, is the lowest realm of rebirth. Of the hell realms, the worst is Avīci (Sanskrit and Pali for "without waves"). The Buddha's disciple, Devadatta, who tried to kill the Buddha on three occasions, as well as create a schism in the monastic order, is said[by whom?] to have been reborn in the Avici Hell.
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+Like all realms of rebirth in Buddhism, rebirth in the Hell realms is not permanent, though suffering can persist for eons before being reborn again.[citation needed] In the Lotus Sutra, the Buddha teaches that eventually even Devadatta will become a Pratyekabuddha himself, emphasizing the temporary nature of the Hell realms. Thus, Buddhism teaches to escape the endless migration of rebirths (both positive and negative) through the attainment of Nirvana.
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+The Bodhisattva Ksitigarbha, according to the Ksitigarbha Sutra, made a great vow as a young girl to not reach Nirvana until all beings were liberated from the Hell Realms or other unwholesome rebirths. In popular literature, Ksitigarbha travels to the Hell realms to teach and relieve beings of their suffering.
+
+Early Vedic religion does not have a concept of Hell. Ṛg-veda mentions three realms, bhūr (the earth), svar (the sky) and bhuvas or antarikṣa (the middle area, i.e. air or atmosphere). In later Hindu literature, especially the law books and Puranas, more realms are mentioned, including a realm similar to Hell, called naraka (in Devanāgarī: नरक). Yama as the first born human (together with his twin sister Yamī), by virtue of precedence, becomes ruler of men and a judge on their departure. Originally he resides in Heaven, but later, especially medieval, traditions mention his court in naraka.[citation needed]
+
+In the law-books (smṛtis and dharma-sūtras, like the Manu-smṛti), naraka is a place of punishment for sins. It is a lower spiritual plane (called naraka-loka) where the spirit is judged and the partial fruits of karma affect the next life. In Mahabharata there is a mention of the Pandavas and the Kauravas both going to Heaven. At first Yudhisthir goes to heaven where he sees Duryodhana enjoying heaven; Indra tells him that Duryodhana is in heaven as he did his Kshatriya duties. Then he shows Yudhisthir hell where it appears his brothers are. Later it is revealed that this was a test for Yudhisthir and that his brothers and the Kauravas are all in heaven and live happily in the divine abode of gods. Hells are also described in various Puranas and other scriptures. The Garuda Purana gives a detailed account of Hell and its features; it lists the amount of punishment for most crimes, much like a modern-day penal code.
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+It is believed[by whom?] that people who commit sins go to Hell and have to go through punishments in accordance with the sins they committed. The god Yamarāja, who is also the god of death, presides over Hell. Detailed accounts of all the sins committed by an individual are kept by Chitragupta, who is the record keeper in Yama's court. Chitragupta reads out the sins committed and Yama orders appropriate punishments to be given to individuals. These punishments include dipping in boiling oil, burning in fire, torture using various weapons, etc. in various Hells. Individuals who finish their quota of the punishments are reborn in accordance with their balance of karma. All created beings are imperfect and thus have at least one sin to their record; but if one has generally led a pious life, one ascends to svarga, a temporary realm of enjoyment similar to Paradise, after a brief period of expiation in Hell and before the next reincarnation, according to the law of karma.[citation needed] With the exception of Hindu philosopher Madhva, time in Hell is not regarded as eternal damnation within Hinduism.[120]
+
+According to Brahma Kumaris, the iron age (Kali Yuga) is regarded as hell.
+
+In Jain cosmology, Naraka (translated as Hell) is the name given to realm of existence having great suffering. However, a Naraka differs from the hells of Abrahamic religions as souls are not sent to Naraka as the result of a divine judgment and punishment. Furthermore, length of a being's stay in a Naraka is not eternal, though it is usually very long and measured in billions of years. A soul is born into a Naraka as a direct result of his or her previous karma (actions of body, speech and mind), and resides there for a finite length of time until his karma has achieved its full result. After his karma is used up, he may be reborn in one of the higher worlds as the result of an earlier karma that had not yet ripened.
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+The Hells are situated in the seven grounds at the lower part of the universe. The seven grounds are:
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+The hellish beings are a type of souls which are residing in these various hells. They are born in hells by sudden manifestation.[121] The hellish beings possess vaikriya  body (protean body which can transform itself and take various forms). They have a fixed life span (ranging from ten thousand to billions of years) in the respective hells where they reside. According to Jain scripture, Tattvarthasutra, following are the causes for birth in hell:[122]
+
+According to Meivazhi, the purpose of all religions is to guide people to Heaven.[124] However, those who do not approach God and are not blessed by Him are believed to be condemned to Hell.[125]
+
+In Sikh thought, Heaven and Hell are not places for living hereafter, they are part of spiritual topography of man and do not exist otherwise. They refer to good and evil stages of life respectively and can be lived now and here during our earthly existence.[126] For example, Guru Arjan explains that people who are entangled in emotional attachment and doubt are living in hell on this Earth i.e. their life is hellish.
+
+So many are being drowned in emotional attachment and doubt; they dwell in the most horrible hell.
+
+Ancient Taoism had no concept of Hell, as morality was seen to be a man-made distinction and there was no concept of an immaterial soul. In its home country China, where Taoism adopted tenets of other religions, popular belief endows Taoist Hell with many deities and spirits who punish sin in a variety of horrible ways.
+
+Diyu is the realm of the dead in Chinese mythology. It is very loosely based upon the Buddhist concept of Naraka combined with traditional Chinese afterlife beliefs and a variety of popular expansions and re-interpretations of these two traditions. Ruled by Yanluo Wang, the King of Hell, Diyu is a maze of underground levels and chambers where souls are taken to atone for their earthly sins.
+
+Incorporating ideas from Taoism and Buddhism as well as traditional Chinese folk religion, Diyu is a kind of purgatory place which serves not only to punish but also to renew spirits ready for their next incarnation. There are many deities associated with the place, whose names and purposes are the subject of much conflicting information.
+
+The exact number of levels in Chinese Hell – and their associated deities – differs according to the Buddhist or Taoist perception. Some speak of three to four 'Courts', other as many as ten. The ten judges are also known as the 10 Kings of Yama. Each Court deals with a different aspect of atonement. For example, murder is punished in one Court, adultery in another. According to some Chinese legends, there are eighteen levels in Hell. Punishment also varies according to belief, but most legends speak of highly imaginative chambers where wrong-doers are sawn in half, beheaded, thrown into pits of filth or forced to climb trees adorned with sharp blades.
+
+However, most legends agree that once a soul (usually referred to as a 'ghost') has atoned for their deeds and repented, he or she is given the Drink of Forgetfulness by Meng Po and sent back into the world to be reborn, possibly as an animal or a poor or sick person, for further punishment.
+
+Zoroastrianism has historically suggested several possible fates for the wicked, including annihilation, purgation in molten metal, and eternal punishment, all of which have standing in Zoroaster's writings. Zoroastrian eschatology includes the belief that wicked souls will remain in Duzakh until, following the arrival of three saviors at thousand-year intervals, Ahura Mazda reconciles the world, destroying evil and resurrecting tormented souls to perfection.[128]
+
+The sacred Gathas mention a "House of the Lie″ for those "that are of an evil dominion, of evil deeds, evil words, evil Self, and evil thought, Liars."[129] However, the best-known Zoroastrian text to describe hell in detail is the Book of Arda Viraf.[130] It depicts particular punishments for particular sins—for instance, being trampled by cattle as punishment for neglecting the needs of work animals.[131] Other descriptions can be found in the Book of Scriptures (Hadhokht Nask), Religious Judgments (Dadestan-i Denig) and the Book of the Judgments of the Spirit of Wisdom (Mainyo-I-Khard).[132]
+
+The Mandaeans believe in purification of souls inside of Leviathan,[133] whom they also call Ur.[134] Within detention houses, so called Mattarathas,[135] the detained souls would receive so much punishment that they would wish to die a Second death, which would, however, not (yet) befall their spirit.[136] At the end of days, the souls of the Mandaeans which could be purified, would be liberated out of Ur's mouth.[137] After this, Ur would get destroyed along with the souls remaining inside him,[138] so they die the second death.[139]
+
+The two oldest sects of Wicca- Gardnerian Wicca and Alexandrian Wicca, include "the wiccan laws" that Gerald Gardner wrote. Those laws state that wiccan souls are privileged with reincarnation, but that the souls of wiccans who break the wiccan laws, "even under torture", would be cursed by the goddess, never be reborn on earth, and "remain where they belong, in the Hell of the Christians."[140][141] Later wiccan sects do not necessarily include Gerald Gardner's wiccan laws. The influential wiccan author Raymond Buckland wrote that the wiccan laws are unimportant. Solitary neo-wiccans, who originated in the 1980s, do not include the wiccan laws in their doctrine.
+
+In his Divina commedia (Divine Comedy), set in the year 1300), Dante Alighieri employed the concept of taking Virgil as his guide through Inferno (and then, in the second canticle, up the mountain of Purgatorio). Virgil himself is not condemned to Hell proper in Dante's poem but is rather, as a virtuous pagan, confined to Limbo just at the edge of Hell. The geography of Hell is very elaborately laid out in this work, with nine concentric rings leading deeper into Earth, and deeper into the various punishments of Hell, until, at the center of the world, Dante finds Satan himself trapped in the frozen lake of Cocytus. A small tunnel leads past Satan and out to the other side of the world, at the base of the Mount of Purgatory.
+
+John Milton's Paradise Lost (1667) opens with the fallen angels, including their leader Satan, waking up in Hell after having been defeated in the war in heaven and the action returns there at several points throughout the poem. Milton portrays Hell as the abode of the demons, and the passive prison from which they plot their revenge upon Heaven through the corruption of the human race. 19th-century French poet Arthur Rimbaud alluded to the concept as well in the title and themes of one of his major works, A Season in Hell. Rimbaud's poetry portrays his own suffering in a poetic form as well as other themes.
+
+Many of the great epics of European literature include episodes that occur in Hell. In the Roman poet Virgil's Latin epic, the Aeneid, Aeneas descends into Dis (the underworld) to visit his father's spirit. The underworld is only vaguely described, with one unexplored path leading to the punishments of Tartarus, while the other leads through Erebus and the Elysian Fields.
+
+The idea of Hell was highly influential to writers such as Jean-Paul Sartre who authored the 1944 play No Exit about the idea that "Hell is other people". Although not a religious man, Sartre was fascinated by his interpretation of a Hellish state of suffering. C.S. Lewis's The Great Divorce (1945) borrows its title from William Blake's Marriage of Heaven and Hell (1793) and its inspiration from the Divine Comedy as the narrator is likewise guided through Hell and Heaven. Hell is portrayed here as an endless, desolate twilight city upon which night is imperceptibly sinking. The night is actually the Apocalypse, and it heralds the arrival of the demons after their judgment. Before the night comes, anyone can escape Hell if they leave behind their former selves and accept Heaven's offer, and a journey to Heaven reveals that Hell is infinitely small; it is nothing more or less than what happens to a soul that turns away from God and into itself.
+
+Piers Anthony in his series Incarnations of Immortality portrays examples of Heaven and Hell via Death, Fate, Underworld, Nature, War, Time, Good-God, and Evil-Devil. Robert A. Heinlein offers a yin-yang version of Hell where there is still some good within; most evident in his book Job: A Comedy of Justice. Lois McMaster Bujold uses her five Gods 'Father, Mother, Son, Daughter and Bastard' in The Curse of Chalion with an example of Hell as formless chaos. Michael Moorcock is one of many who offer Chaos-Evil-(Hell) and Uniformity-Good-(Heaven) as equally unacceptable extremes which must be held in balance; in particular in the Elric and Eternal Champion series. Fredric Brown wrote a number of fantasy short stories about Satan's activities in Hell. Cartoonist Jimmy Hatlo created a series of cartoons about life in Hell called The Hatlo Inferno, which ran from 1953 to 1958.[142]

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+In religion and folklore, Hell is an afterlife location in which evil souls are subjected to punitive suffering, often torture as eternal punishment after death. Religions with a linear divine history often depict hells as eternal destinations, the biggest examples of which are Christianity and Islam, whereas religions with reincarnation usually depict a hell as an intermediary period between incarnations, as is the case in the dharmic religions. Religions typically locate hell in another dimension or under Earth's surface. Other afterlife destinations include Heaven, Paradise, Purgatory, Limbo, and the underworld.
+
+Other religions, which do not conceive of the afterlife as a place of punishment or reward, merely describe an abode of the dead, the grave, a neutral place that is located under the surface of Earth (for example, see Kur, Hades, and Sheol). Such places are sometimes equated with the English word hell, though a more correct translation would be "underworld" or "world of the dead". The ancient Mesopotamian, Greek, Roman, and Finnic religions include entrances to the underworld from the land of the living.
+
+The modern English word hell is derived from Old English hel, helle (first attested around 725 AD to refer to a nether world of the dead) reaching into the Anglo-Saxon pagan period.[1] The word has cognates in all branches of the Germanic languages, including Old Norse hel (which refers to both a location and goddess-like being in Norse mythology), Old Frisian helle, Old Saxon hellia, Old High German hella, and Gothic halja. All forms ultimately derive from the reconstructed Proto-Germanic feminine noun *xaljō or *haljō ('concealed place, the underworld'). In turn, the Proto-Germanic form derives from the o-grade form of the Proto-Indo-European root *kel-, *kol-: 'to cover, conceal, save'.[2] Indo-European cognates include Latin cēlāre ("to hide", related to the English word cellar) and early Irish ceilid ("hides"). Upon the Christianization of the Germanic peoples, extension of Proto-Germanic *xaljō were reinterpreted to denote the underworld in Christian mythology,[1][3] for which see Gehenna.
+
+Related early Germanic terms and concepts include Proto-Germanic *xalja-rūnō(n), a feminine compound noun, and *xalja-wītjan, a neutral compound noun. This form is reconstructed from the Latinized Gothic plural noun *haliurunnae (attested by Jordanes; according to philologist Vladimir Orel, meaning 'witches'), Old English helle-rúne ('sorceress, necromancer', according to Orel), and Old High German helli-rūna 'magic'. The compound is composed of two elements: *xaljō (*haljō) and *rūnō, the Proto-Germanic precursor to Modern English rune.[4] The second element in the Gothic haliurunnae may however instead be an agent noun from the verb rinnan ("to run, go"), which would make its literal meaning "one who travels to the netherworld".[5][6]
+
+Proto-Germanic *xalja-wītjan (or *halja-wītjan) is reconstructed from Old Norse hel-víti 'hell', Old English helle-wíte 'hell-torment, hell', Old Saxon helli-wīti 'hell', and the Middle High German feminine noun helle-wīze. The compound is a compound of *xaljō (discussed above) and *wītjan (reconstructed from forms such as Old English witt 'right mind, wits', Old Saxon gewit 'understanding', and Gothic un-witi 'foolishness, understanding').[7]
+
+Hell appears in several mythologies and religions. It is commonly inhabited by demons and the souls of dead people. A fable about Hell which recurs in folklore across several cultures is the allegory of the long spoons. Hell is often depicted in art and literature, perhaps most famously in Dante's Divine Comedy.
+
+Punishment in Hell typically corresponds to sins committed during life. Sometimes these distinctions are specific, with damned souls suffering for each sin committed (see for example Plato's myth of Er or Dante's The Divine Comedy), but sometimes they are general, with condemned sinners relegated to one or more chamber of Hell or to a level of suffering.
+
+In many religious cultures, including Christianity and Islam, Hell is often depicted as fiery, painful, and harsh, inflicting suffering on the guilty. Despite these common depictions of Hell as a place of fire, some other traditions portray Hell as cold. Buddhist – and particularly Tibetan Buddhist – descriptions of Hell feature an equal number of hot and cold Hells. Among Christian descriptions Dante's Inferno portrays the innermost (9th) circle of Hell as a frozen lake of blood and guilt.[11]
+But cold also played a part in earlier Christian depictions of Hell, beginning with the Apocalypse of Paul, originally from the early third century;[12] the "Vision of Dryhthelm" by the Venerable Bede from the seventh century;[13] "St Patrick's Purgatory", "The Vision of Tundale" or "Visio Tnugdali", and the "Vision of the Monk of Eynsham", all from the twelfth century;[14] and the "Vision of Thurkill" from the early thirteenth century.[15]
+
+The Sumerian afterlife was a dark, dreary cavern located deep below the ground,[16] where inhabitants were believed to continue "a shadowy version of life on earth".[16] This bleak domain was known as Kur,[17]:114 and was believed to be ruled by the goddess Ereshkigal.[16][18]:184 All souls went to the same afterlife,[16] and a person's actions during life had no effect on how the person would be treated in the world to come.[16]
+
+The souls in Kur were believed to eat nothing but dry dust[17]:58 and family members of the deceased would ritually pour libations into the dead person's grave through a clay pipe, thereby allowing the dead to drink.[17]:58 Nonetheless, funerary evidence indicates that some people believed that the goddess Inanna, Ereshkigal's younger sister, had the power to award her devotees with special favors in the afterlife.[16][19] During the Third Dynasty of Ur, it was believed that a person's treatment in the afterlife depended on how he or she was buried;[17]:58 those that had been given sumptuous burials would be treated well,[17]:58 but those who had been given poor burials would fare poorly.[17]:58
+
+The entrance to Kur was believed to be located in the Zagros mountains in the far east.[17]:114 It had seven gates, through which a soul needed to pass.[16] The god Neti was the gatekeeper.[18]:184[17]:86 Ereshkigal's sukkal, or messenger, was the god Namtar.[17]:134[18]:184 Galla were a class of demons that were believed to reside in the underworld;[17]:85 their primary purpose appears to have been to drag unfortunate mortals back to Kur.[17]:85 They are frequently referenced in magical texts,[17]:85–86 and some texts describe them as being seven in number.[17]:85–86 Several extant poems describe the galla dragging the god Dumuzid into the underworld.[17]:86 The later Mesopotamians knew this underworld by its East Semitic name: Irkalla. During the Akkadian Period, Ereshkigal's role as the ruler of the underworld was assigned to Nergal, the god of death.[16][18]:184 The Akkadians attempted to harmonize this dual rulership of the underworld by making Nergal Ereshkigal's husband.[16]
+
+With the rise of the cult of Osiris during the Middle Kingdom the "democratization of religion" offered to even his humblest followers the prospect of eternal life, with moral fitness becoming the dominant factor in determining a person's suitability. At death a person faced judgment by a tribunal of forty-two divine judges. If they had led a life in conformance with the precepts of the goddess Maat, who represented truth and right living, the person was welcomed into the heavenly reed fields. If found guilty the person was thrown to Ammit, the "devourer of the dead" and would be condemned to the lake of fire.[21] The person taken by the devourer is subject first to terrifying punishment and then annihilated. These depictions of punishment may have influenced medieval perceptions of the inferno in hell via early Christian and Coptic texts.[22] Purification for those considered justified appears in the descriptions of "Flame Island", where humans experience the triumph over evil and rebirth. For the damned complete destruction into a state of non-being awaits but there is no suggestion of eternal torture; the weighing of the heart in Egyptian mythology can lead to annihilation.[23][24] The Tale of Khaemwese describes the torment of a rich man, who lacked charity, when he dies and compares it to the blessed state of a poor man who has also died.[25]
+Divine pardon at judgement always remained a central concern for the ancient Egyptians.[26]
+
+Modern understanding of Egyptian notions of hell relies on six ancient texts:[27]
+
+In classic Greek mythology, below Heaven, Earth, and Pontus is Tartarus, or Tartaros (Greek Τάρταρος, deep place). It is either a deep, gloomy place, a pit or abyss used as a dungeon of torment and suffering that resides within Hades (the entire underworld) with Tartarus being the hellish component. In the Gorgias, Plato (c. 400 BC) wrote that souls of the deceased were judged after they payed for crossing the river of the dead and those who received punishment were sent to Tartarus.[28] As a place of punishment, it can be considered a hell. The classic Hades, on the other hand, is more similar to Old Testament Sheol. The Romans later adopted these views.
+
+The hells of Europe include Breton mythology's "Anaon", Celtic mythology's "Uffern", Slavic mythology's "Peklo", the hell of Sami mythology and Finnish "tuonela" ("manala").
+
+The hells of Asia include the Bagobo "Gimokodan" (which is believed to be more of an otherworld, where the Red Region is reserved who those who died in battle, while ordinary people go to the White Region)[29] and ancient Indian mythology's "Kalichi" or "Naraka".
+
+According to a few sources, hell is below ground, and described as an uninviting wet[30] or fiery place reserved for sinful people in the Ainu religion, as stated by missionary John Batchelor.[31] However, belief in hell does not appear in oral tradition of the Ainu.[32] Instead, there is belief within the Ainu religion that the soul of the deceased (ramat) would become a kamuy after death.[32] There is also belief that the soul of someone who has been wicked during lifetime, committed suicide, got murdered or died in great agony would become a ghost (tukap) who would haunt the living,[32] to come to fulfillment from which it was excluded during life.[33]
+
+In Tengrism, it was believed that the wicked would get punished in Tamag before they would be brought to the third floor of the sky.[34]
+
+In Taoism, hell is represented by Diyu.
+
+The Hell of Swahili mythology is called kuzimu, and believe in it deleveloped in the 7th and 8th century under the influence of Muslim merchants at the east African coast.[35] It is imagined as a very cold place.[35] Serer religion rejects the general notion of heaven and hell.[36] In Serer religion, acceptance by the ancestors who have long departed is as close to any heaven as one can get. Rejection and becoming a wandering soul is a sort of hell for one passing over. The souls of the dead must make their way to Jaaniw (the sacred dwelling place of the soul). Only those who have lived their lives on earth in accordance with Serer doctrines will be able to make this necessary journey and thus accepted by the ancestors. Those who can't make the journey become lost and wandering souls, but they do not burn in "hell fire".[36][37]
+
+The hells of the Americas include the Aztec religion's Mictlan, Inuit religion's Adlivun, and the Yanomami religion's Shobari Waka. In Mayan religion, Xibalba (or Metnal) is the dangerous underworld of nine levels. The road into and out of it is said to be steep, thorny and very forbidding. Ritual healers would intone healing prayers banishing diseases to Xibalba. Much of the Popol Vuh describes the adventures of the Maya Hero Twins in their cunning struggle with the evil lords of Xibalba.
+
+The Aztecs believed that the dead traveled to Mictlan, a neutral place found far to the north. There was also a legend of a place of white flowers, which was always dark, and was home to the gods of death, particularly Mictlantecutli and his spouse Mictlantecihuatl, which means literally "lords of Mictlan". The journey to Mictlan took four years, and the travelers had to overcome difficult tests, such as passing a mountain range where the mountains crashed into each other, a field where the wind carried flesh-scraping knives, and a river of blood with fearsome jaguars.
+
+In pre-Christian Fijian mythology there was belief in an underworld called Murimuria.
+
+Hell is conceived of in most Abrahamic religions as a place of, or a form of, punishment.[38]
+
+Early Judaism had no concept of Hell, although the concept of an afterlife was introduced during the Hellenistic period, apparently from neighboring Hellenistic religions. It occurs for example in the Book of Daniel. Daniel 12:2 proclaims "And many of those who sleep in the dust of the earth shall awake, Some to everlasting life, Some to shame and everlasting contempt."
+
+Judaism does not have a specific doctrine about the afterlife, but it does have a mystical/Orthodox tradition of describing Gehinnom. Gehinnom is not Hell, but originally a grave and in later times a sort of Purgatory where one is judged based on one's life's deeds, or rather, where one becomes fully aware of one's own shortcomings and negative actions during one's life. The Kabbalah explains it as a "waiting room" (commonly translated as an "entry way") for all souls (not just the wicked). The overwhelming majority of rabbinic thought maintains that people are not in Gehinnom forever; the longest that one can be there is said to be 12 months, however there has been the occasional noted exception. Some consider it a spiritual forge where the soul is purified for its eventual ascent to Olam Habah (heb. עולם הבא; lit. "The world to come", often viewed as analogous to heaven). This is also mentioned in the Kabbalah, where the soul is described as breaking, like the flame of a candle lighting another: the part of the soul that ascends being pure and the "unfinished" piece being reborn.
+
+According to Jewish teachings, hell is not entirely physical; rather, it can be compared to a very intense feeling of shame. People are ashamed of their misdeeds and this constitutes suffering which makes up for the bad deeds. When one has so deviated from the will of God, one is said to be in Gehinnom. This is not meant to refer to some point in the future, but to the very present moment. The gates of teshuva (return) are said to be always open, and so one can align his will with that of God at any moment. Being out of alignment with God's will is itself a punishment according to the Torah.
+
+Many scholars of Jewish mysticism, particularly of the Kabbalah, describe seven "compartments" or "habitations" of Hell, just as they describe seven divisions of Heaven. These divisions go by many different names, and the most frequently mentioned are as follows:[39]
+
+Besides those mentioned above, there also exist additional terms that have been often used to either refer to Hell in general or to some region of the underworld:
+
+For more information, see Qliphoth.
+
+Maimonides declares in his 13 principles of faith that the hells of the rabbinic literature were pedagocically motivated inventions to encourage respect of the Torah commandements by mankind, which had been regarded as immature.[49] Instead of being sent to hell, the souls of the wicked would actually get annihilated.[50]
+
+The Christian doctrine of hell derives from passages in the New Testament. The word hell does not appear in the Greek New Testament; instead one of three words is used: the Greek words Tartarus or Hades, or the Hebrew word Gehinnom.
+
+In the Septuagint and New Testament the authors used the Greek term Hades for the Hebrew Sheol, but often with Jewish rather than Greek concepts in mind. In the Jewish concept of Sheol, such as expressed in Ecclesiastes,[51] Sheol or Hades is a place where there is no activity. However, since Augustine, some[which?] Christians have believed that the souls of those who die either rest peacefully, in the case of Christians, or are afflicted, in the case of the damned, after death until the resurrection.[52]
+
+While these three terms are translated in the KJV as "hell" these three terms have three very different meanings.
+
+The Roman Catholic Church defines Hell as "a state of definitive self-exclusion from communion with God and the blessed." One finds oneself in Hell as the result of dying in mortal sin without repenting and accepting God's merciful love, becoming eternally separated from him by one's own free choice[66] immediately after death.[67] In the Roman Catholic Church, many other Christian churches, such as the Baptists and Episcopalians, and some Greek Orthodox churches,[68] Hell is taught as the final destiny of those who have not been found worthy after the general resurrection and last judgment,[69][70][71] where they will be eternally punished for sin and permanently separated from God. The nature of this judgment is inconsistent with many Protestant churches teaching the saving comes from accepting Jesus Christ as their savior, while the Greek Orthodox and Catholic Churches teach that the judgment hinges on both faith and works. However, many Liberal Christians throughout Liberal Protestant and Anglican churches believe in universal reconciliation (see below), even though it contradicts the traditional doctrines that are usually held by the evangelicals within their denominations.[72]
+
+Some modern Christian theologians subscribe to the doctrines of conditional immortality. Conditional immortality is the belief that the soul dies with the body and does not live again until the resurrection. As with other Jewish writings of the Second Temple period, the New Testament text distinguishes two words, both translated "Hell" in older English Bibles: Hades, "the grave", and Gehenna where God "can destroy both body and soul".[73] A minority of Christians read this to mean that neither Hades nor Gehenna are eternal but refer to the ultimate destruction of the wicked in the Lake of Fire in a consuming fire after resurrection. However, because of the Greek words used in translating from the Hebrew text, the Hebrew ideas have become confused with Greek myths and ideas.  In the Hebrew text when people died they went to Sheol, the grave[74] and the wicked ultimately went to Gehenna and were consumed by fire. The Hebrew words for "the grave" or "death" or "eventual destruction of the wicked", were translated using Greek words and later texts became a mix of mistranslation, pagan influence, and Greek myth.[75]
+
+Christian mortalism is the doctrine that all men and women, including Christians, must die, and do not continue and are not conscious after death. Therefore, annihilationism includes the doctrine that "the wicked" are also destroyed rather than tormented forever in traditional "Hell" or the lake of fire. Christian mortalism and annihilationism are directly related to the doctrine of conditional immortality, the idea that a human soul is not immortal unless it is given eternal life at the second coming of Christ and resurrection of the dead.
+
+Biblical scholars looking at the issue through the Hebrew text have denied the teaching of innate immortality.[76][77] Rejection of the immortality of the soul, and advocacy of Christian mortalism, was a feature of Protestantism since the early days of the Reformation with Martin Luther himself rejecting the traditional idea, though his mortalism did not carry into orthodox Lutheranism. One of the most notable English opponents of the immortality of the soul was Thomas Hobbes who describes the idea as a Greek "contagion" in Christian doctrine.[78] Modern proponents of conditional immortality include some in the Anglican church such as N.T. Wright[79] and as denominations the Seventh-day Adventists, Bible Students, Jehovah's Witnesses, Christadelphians, Living Church of God, The Church of God International, and some other Protestant Christians, as well as recent Roman Catholic teaching. It is not Roman Catholic dogma that anyone is in Hell,[80] though many individual Catholics do not share this view. The 1993 Catechism of the Catholic Church states:[81] "This state of definitive self-exclusion from communion with God and the blessed is called 'hell'" and[82] "... they suffer the punishments of hell, "eternal fire." The chief punishment of hell is eternal separation from God" (CCC 1035).  During an Audience in 1999, Pope John Paul II commented: "images of hell that Sacred Scripture presents to us must be correctly interpreted. They show the complete frustration and emptiness of life without God. Rather than a place, hell indicates the state of those who freely and definitively separate themselves from God, the source of all life and joy."[83]
+
+The Seventh-day Adventist Church's official beliefs support annihilationism.[84][85] They deny the Catholic purgatory and teach that the dead lie in the grave until they are raised for a last judgment, both the righteous and wicked await the resurrection at the Second Coming. Seventh-day Adventists believe that death is a state of unconscious sleep until the resurrection. They base this belief on biblical texts such as Ecclesiastes 9:5 which states "the dead know nothing", and 1 Thessalonians 4:13–18 which contains a description of the dead being raised from the grave at the second coming. These verses, it is argued, indicate that death is only a period or form of slumber.
+
+Adventists teach that the resurrection of the righteous will take place shortly after the second coming of Jesus, as described in Revelation 20:4–6 that follows Revelation 19:11–16, whereas the resurrection of the wicked will occur after the millennium, as described in Revelation 20:5 and 20:12–13 that follow Revelation 20:4 and 6–7, though Revelation 20:12–13 and 15 actually describe a mixture of saved and condemned people being raised from the dead and judged. Adventists reject the traditional doctrine of hell as a state of everlasting conscious torment, believing instead that the wicked will be permanently destroyed after the millennium by the lake of fire, which is called 'the second death' in Revelation 20:14.
+
+Those Adventist doctrines about death and hell reflect an underlying belief in: (a) conditional immortality (or conditionalism), as opposed to the immortality of the soul; and (b) the monistic nature of human beings, in which the soul is not separable from the body, as opposed to bipartite or tripartite conceptions, in which the soul is separable.
+
+Jehovah's Witnesses hold that the soul ceases to exist when the person dies[86] and therefore that Hell (Sheol or Hades) is a state of non-existence.[86] In their theology, Gehenna differs from Sheol or Hades in that it holds no hope of a resurrection.[86] Tartarus is held to be the metaphorical state of debasement of the fallen angels between the time of their moral fall (Genesis chapter 6) until their post-millennial destruction along with Satan (Revelation chapter 20).[87]
+
+Bible Students and Christadelphians also believe in annihilationism.
+
+Christian Universalists believe in universal reconciliation, the belief that all human souls will be eventually reconciled with God and admitted to Heaven.[88] This belief is held by some Unitarian-Universalists.[89][90][91]
+
+According to Emanuel Swedenborg's Second Coming Christian revelation, hell exists because evil people want it.[92] They, not God, introduced evil to the human race.[93] In Swedenborgianism, every soul joins the like-minded group after death in which it feels the most comfortable. Hell is therefore believed to be a place of hapiness for the souls which delight in evilness.[94]
+
+Members of The Church of Jesus Christ of Latter-day Saints (LDS Church) teach that hell is a state between death and resurrection, in which those spirits who did not repent while on earth must suffer for their own sins (Doctrine and Covenants 19:15–17[95]). After that, only the Sons of perdition, who committed the Eternal sin, would be cast into Outer darkness. However, according to Mormon faith, committing the Eternal sin requires so much knowledge that most persons cannot do this.[96] Satan and Cain are counted as examples of Sons of perdition.
+
+In Islam, jahannam (in Arabic: جهنم) (related to the Hebrew word gehinnom) is the counterpart to heaven and likewise divided into seven layers, both co-existing with the temporal world,[97] filled with blazing fire, boiling water, and a variety of other torments for those who have been condemned to it in the hereafter. In the Quran, God declares that the fire of Jahannam is prepared for both mankind and jinn.[98][99] After the Day of Judgement, it is to be occupied by those who do not believe in God, those who have disobeyed his laws, or rejected his messengers.[100] "Enemies of Islam" are sent to Hell immediately upon their deaths.[101] Muslim modernists downplay the vivid descriptions of hell common during Classical period, on one hand reaffirming that the afterlife must not be denied, but simultaneously asserting its exact nature remains unknown. Other modern Muslims continue the line of Sufism as an interiorized hell, combining the eschatological thoughts of Ibn Arabi and Rumi with Western philosophy.[102] Although disputed by some scholars, most scholars consider jahannam to be eternal.[103][104] There is belief that the fire which represents the own bad deeds can already be seen during the Punishment of the Grave, and that the spiritual pain caused by this can lead to purification of the soul.[105]
+
+Medieval sources usually identified hell with the seven layers of the earth mentioned in Surah 65:12, inhabited by devils, harsh angels, scorpions and serpents, who torment the sinners. They described thorny shrubs, seas filled with blood and fire and darkness only illuminated by the flames of hell.[106] However, some sources also mention a place of extreme cold at the bottom of hell, called Zamhareer, characterized in as being unbearably cold, with blizzards, ice, and snow.[107] Maalik is thought of as the keeper of the gates of hell, namely appears in Ibn Abbas' Isra and Mi'raj.[108] Over hell, a narrow bridge called As-Sirāt is spanned. On Judgement Day one must pass over it to reach paradise, but those destined for hell will find too narrow and fall from into their new abode.[109] Iblis, the temporary ruler of hell,[110] is thought of residing in the bottom of hell, from where he commands his hosts of infernal demons.[111][112] But contrary to Christian traditions, Iblis and his infernal hosts do not wage war against God,[113] his enmity applies against humanity only. Further, his dominion in hell is also his punishment. According to the Muwatta Hadith, the Bukhari Hadith, the Tirmidhi Hadith, and the Kabir Hadith, Muhammad claimed that the fire of Jahannam is not red, but pitch-black, and is 70 times hotter than ordinary fire, and is much more painful than ordinary fire.[citation needed]
+
+Polytheism (shirk) is regarded as a particularly grievous sin; therefore entering Paradise is forbidden to a polytheist (musyrik) because his place is Hell;[114] and the lowest pit of Hell (Hawiyah), is intended for hypocrites who claimed aloud to believe in God and his messenger but in their hearts did not.[115] Not all Muslims and scholars agree whether hell is an eternal destination or whether some or all of the condemned will eventually be forgiven and allowed to enter paradise.[101][113][116][117]
+
+In the Bahá'í Faith, the conventional descriptions of Hell and Heaven are considered to be symbolic representations of spiritual conditions. The Bahá'í writings describe closeness to God to be Heaven, and conversely, remoteness from God as Hell.[118] The Bahá'í writings state that the soul is immortal and after death it will continue to progress until it finally attains God's presence.[119]
+
+In "Devaduta Sutta", the 130th discourse of the Majjhima Nikaya, Buddha teaches about hell in vivid detail. Buddhism teaches that there are five[citation needed] (sometimes six[citation needed]) realms of rebirth, which can then be further subdivided into degrees of agony or pleasure. Of these realms, the hell realms, or Naraka, is the lowest realm of rebirth. Of the hell realms, the worst is Avīci (Sanskrit and Pali for "without waves"). The Buddha's disciple, Devadatta, who tried to kill the Buddha on three occasions, as well as create a schism in the monastic order, is said[by whom?] to have been reborn in the Avici Hell.
+
+Like all realms of rebirth in Buddhism, rebirth in the Hell realms is not permanent, though suffering can persist for eons before being reborn again.[citation needed] In the Lotus Sutra, the Buddha teaches that eventually even Devadatta will become a Pratyekabuddha himself, emphasizing the temporary nature of the Hell realms. Thus, Buddhism teaches to escape the endless migration of rebirths (both positive and negative) through the attainment of Nirvana.
+
+The Bodhisattva Ksitigarbha, according to the Ksitigarbha Sutra, made a great vow as a young girl to not reach Nirvana until all beings were liberated from the Hell Realms or other unwholesome rebirths. In popular literature, Ksitigarbha travels to the Hell realms to teach and relieve beings of their suffering.
+
+Early Vedic religion does not have a concept of Hell. Ṛg-veda mentions three realms, bhūr (the earth), svar (the sky) and bhuvas or antarikṣa (the middle area, i.e. air or atmosphere). In later Hindu literature, especially the law books and Puranas, more realms are mentioned, including a realm similar to Hell, called naraka (in Devanāgarī: नरक). Yama as the first born human (together with his twin sister Yamī), by virtue of precedence, becomes ruler of men and a judge on their departure. Originally he resides in Heaven, but later, especially medieval, traditions mention his court in naraka.[citation needed]
+
+In the law-books (smṛtis and dharma-sūtras, like the Manu-smṛti), naraka is a place of punishment for sins. It is a lower spiritual plane (called naraka-loka) where the spirit is judged and the partial fruits of karma affect the next life. In Mahabharata there is a mention of the Pandavas and the Kauravas both going to Heaven. At first Yudhisthir goes to heaven where he sees Duryodhana enjoying heaven; Indra tells him that Duryodhana is in heaven as he did his Kshatriya duties. Then he shows Yudhisthir hell where it appears his brothers are. Later it is revealed that this was a test for Yudhisthir and that his brothers and the Kauravas are all in heaven and live happily in the divine abode of gods. Hells are also described in various Puranas and other scriptures. The Garuda Purana gives a detailed account of Hell and its features; it lists the amount of punishment for most crimes, much like a modern-day penal code.
+
+It is believed[by whom?] that people who commit sins go to Hell and have to go through punishments in accordance with the sins they committed. The god Yamarāja, who is also the god of death, presides over Hell. Detailed accounts of all the sins committed by an individual are kept by Chitragupta, who is the record keeper in Yama's court. Chitragupta reads out the sins committed and Yama orders appropriate punishments to be given to individuals. These punishments include dipping in boiling oil, burning in fire, torture using various weapons, etc. in various Hells. Individuals who finish their quota of the punishments are reborn in accordance with their balance of karma. All created beings are imperfect and thus have at least one sin to their record; but if one has generally led a pious life, one ascends to svarga, a temporary realm of enjoyment similar to Paradise, after a brief period of expiation in Hell and before the next reincarnation, according to the law of karma.[citation needed] With the exception of Hindu philosopher Madhva, time in Hell is not regarded as eternal damnation within Hinduism.[120]
+
+According to Brahma Kumaris, the iron age (Kali Yuga) is regarded as hell.
+
+In Jain cosmology, Naraka (translated as Hell) is the name given to realm of existence having great suffering. However, a Naraka differs from the hells of Abrahamic religions as souls are not sent to Naraka as the result of a divine judgment and punishment. Furthermore, length of a being's stay in a Naraka is not eternal, though it is usually very long and measured in billions of years. A soul is born into a Naraka as a direct result of his or her previous karma (actions of body, speech and mind), and resides there for a finite length of time until his karma has achieved its full result. After his karma is used up, he may be reborn in one of the higher worlds as the result of an earlier karma that had not yet ripened.
+
+The Hells are situated in the seven grounds at the lower part of the universe. The seven grounds are:
+
+The hellish beings are a type of souls which are residing in these various hells. They are born in hells by sudden manifestation.[121] The hellish beings possess vaikriya  body (protean body which can transform itself and take various forms). They have a fixed life span (ranging from ten thousand to billions of years) in the respective hells where they reside. According to Jain scripture, Tattvarthasutra, following are the causes for birth in hell:[122]
+
+According to Meivazhi, the purpose of all religions is to guide people to Heaven.[124] However, those who do not approach God and are not blessed by Him are believed to be condemned to Hell.[125]
+
+In Sikh thought, Heaven and Hell are not places for living hereafter, they are part of spiritual topography of man and do not exist otherwise. They refer to good and evil stages of life respectively and can be lived now and here during our earthly existence.[126] For example, Guru Arjan explains that people who are entangled in emotional attachment and doubt are living in hell on this Earth i.e. their life is hellish.
+
+So many are being drowned in emotional attachment and doubt; they dwell in the most horrible hell.
+
+Ancient Taoism had no concept of Hell, as morality was seen to be a man-made distinction and there was no concept of an immaterial soul. In its home country China, where Taoism adopted tenets of other religions, popular belief endows Taoist Hell with many deities and spirits who punish sin in a variety of horrible ways.
+
+Diyu is the realm of the dead in Chinese mythology. It is very loosely based upon the Buddhist concept of Naraka combined with traditional Chinese afterlife beliefs and a variety of popular expansions and re-interpretations of these two traditions. Ruled by Yanluo Wang, the King of Hell, Diyu is a maze of underground levels and chambers where souls are taken to atone for their earthly sins.
+
+Incorporating ideas from Taoism and Buddhism as well as traditional Chinese folk religion, Diyu is a kind of purgatory place which serves not only to punish but also to renew spirits ready for their next incarnation. There are many deities associated with the place, whose names and purposes are the subject of much conflicting information.
+
+The exact number of levels in Chinese Hell – and their associated deities – differs according to the Buddhist or Taoist perception. Some speak of three to four 'Courts', other as many as ten. The ten judges are also known as the 10 Kings of Yama. Each Court deals with a different aspect of atonement. For example, murder is punished in one Court, adultery in another. According to some Chinese legends, there are eighteen levels in Hell. Punishment also varies according to belief, but most legends speak of highly imaginative chambers where wrong-doers are sawn in half, beheaded, thrown into pits of filth or forced to climb trees adorned with sharp blades.
+
+However, most legends agree that once a soul (usually referred to as a 'ghost') has atoned for their deeds and repented, he or she is given the Drink of Forgetfulness by Meng Po and sent back into the world to be reborn, possibly as an animal or a poor or sick person, for further punishment.
+
+Zoroastrianism has historically suggested several possible fates for the wicked, including annihilation, purgation in molten metal, and eternal punishment, all of which have standing in Zoroaster's writings. Zoroastrian eschatology includes the belief that wicked souls will remain in Duzakh until, following the arrival of three saviors at thousand-year intervals, Ahura Mazda reconciles the world, destroying evil and resurrecting tormented souls to perfection.[128]
+
+The sacred Gathas mention a "House of the Lie″ for those "that are of an evil dominion, of evil deeds, evil words, evil Self, and evil thought, Liars."[129] However, the best-known Zoroastrian text to describe hell in detail is the Book of Arda Viraf.[130] It depicts particular punishments for particular sins—for instance, being trampled by cattle as punishment for neglecting the needs of work animals.[131] Other descriptions can be found in the Book of Scriptures (Hadhokht Nask), Religious Judgments (Dadestan-i Denig) and the Book of the Judgments of the Spirit of Wisdom (Mainyo-I-Khard).[132]
+
+The Mandaeans believe in purification of souls inside of Leviathan,[133] whom they also call Ur.[134] Within detention houses, so called Mattarathas,[135] the detained souls would receive so much punishment that they would wish to die a Second death, which would, however, not (yet) befall their spirit.[136] At the end of days, the souls of the Mandaeans which could be purified, would be liberated out of Ur's mouth.[137] After this, Ur would get destroyed along with the souls remaining inside him,[138] so they die the second death.[139]
+
+The two oldest sects of Wicca- Gardnerian Wicca and Alexandrian Wicca, include "the wiccan laws" that Gerald Gardner wrote. Those laws state that wiccan souls are privileged with reincarnation, but that the souls of wiccans who break the wiccan laws, "even under torture", would be cursed by the goddess, never be reborn on earth, and "remain where they belong, in the Hell of the Christians."[140][141] Later wiccan sects do not necessarily include Gerald Gardner's wiccan laws. The influential wiccan author Raymond Buckland wrote that the wiccan laws are unimportant. Solitary neo-wiccans, who originated in the 1980s, do not include the wiccan laws in their doctrine.
+
+In his Divina commedia (Divine Comedy), set in the year 1300), Dante Alighieri employed the concept of taking Virgil as his guide through Inferno (and then, in the second canticle, up the mountain of Purgatorio). Virgil himself is not condemned to Hell proper in Dante's poem but is rather, as a virtuous pagan, confined to Limbo just at the edge of Hell. The geography of Hell is very elaborately laid out in this work, with nine concentric rings leading deeper into Earth, and deeper into the various punishments of Hell, until, at the center of the world, Dante finds Satan himself trapped in the frozen lake of Cocytus. A small tunnel leads past Satan and out to the other side of the world, at the base of the Mount of Purgatory.
+
+John Milton's Paradise Lost (1667) opens with the fallen angels, including their leader Satan, waking up in Hell after having been defeated in the war in heaven and the action returns there at several points throughout the poem. Milton portrays Hell as the abode of the demons, and the passive prison from which they plot their revenge upon Heaven through the corruption of the human race. 19th-century French poet Arthur Rimbaud alluded to the concept as well in the title and themes of one of his major works, A Season in Hell. Rimbaud's poetry portrays his own suffering in a poetic form as well as other themes.
+
+Many of the great epics of European literature include episodes that occur in Hell. In the Roman poet Virgil's Latin epic, the Aeneid, Aeneas descends into Dis (the underworld) to visit his father's spirit. The underworld is only vaguely described, with one unexplored path leading to the punishments of Tartarus, while the other leads through Erebus and the Elysian Fields.
+
+The idea of Hell was highly influential to writers such as Jean-Paul Sartre who authored the 1944 play No Exit about the idea that "Hell is other people". Although not a religious man, Sartre was fascinated by his interpretation of a Hellish state of suffering. C.S. Lewis's The Great Divorce (1945) borrows its title from William Blake's Marriage of Heaven and Hell (1793) and its inspiration from the Divine Comedy as the narrator is likewise guided through Hell and Heaven. Hell is portrayed here as an endless, desolate twilight city upon which night is imperceptibly sinking. The night is actually the Apocalypse, and it heralds the arrival of the demons after their judgment. Before the night comes, anyone can escape Hell if they leave behind their former selves and accept Heaven's offer, and a journey to Heaven reveals that Hell is infinitely small; it is nothing more or less than what happens to a soul that turns away from God and into itself.
+
+Piers Anthony in his series Incarnations of Immortality portrays examples of Heaven and Hell via Death, Fate, Underworld, Nature, War, Time, Good-God, and Evil-Devil. Robert A. Heinlein offers a yin-yang version of Hell where there is still some good within; most evident in his book Job: A Comedy of Justice. Lois McMaster Bujold uses her five Gods 'Father, Mother, Son, Daughter and Bastard' in The Curse of Chalion with an example of Hell as formless chaos. Michael Moorcock is one of many who offer Chaos-Evil-(Hell) and Uniformity-Good-(Heaven) as equally unacceptable extremes which must be held in balance; in particular in the Elric and Eternal Champion series. Fredric Brown wrote a number of fantasy short stories about Satan's activities in Hell. Cartoonist Jimmy Hatlo created a series of cartoons about life in Hell called The Hatlo Inferno, which ran from 1953 to 1958.[142]

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+English is a West Germanic language that was first spoken in early medieval England and eventually became a global lingua franca.[4][5] It is named after the Angles, one of the ancient Germanic peoples that migrated to the area of Great Britain that later took their name, England. Both names derive from Anglia, a peninsula on the Baltic Sea. English is most closely related to Frisian and Low Saxon, while its vocabulary has been significantly influenced by other Germanic languages, particularly Old Norse (a North Germanic language), as well as Latin and French.[6][7][8]
+
+English has developed over the course of more than 1,400 years. The earliest forms of English, a group of West Germanic (Ingvaeonic) dialects brought to Great Britain by Anglo-Saxon settlers in the 5th century, are collectively called Old English. Middle English began in the late 11th century with the Norman conquest of England; this was a period in which English was influenced by Old French, in particular through its Old Norman dialect.[9][10] Early Modern English began in the late 15th century with the introduction of the printing press to London, the printing of the King James Bible and the start of the Great Vowel Shift.[11]
+
+Modern English has been spreading around the world since the 17th century by the worldwide influence of the British Empire and the United States. Through all types of printed and electronic media of these countries, English has become the leading language of international discourse and the lingua franca in many regions and professional contexts such as science, navigation and law.[12] Modern English grammar is the result of a gradual change from a typical Indo-European dependent marking pattern, with a rich inflectional morphology and relatively free word order, to a mostly analytic pattern with little inflection, a fairly fixed subject–verb–object word order and a complex syntax.[13] Modern English relies more on auxiliary verbs and word order for the expression of complex tenses, aspect and mood, as well as passive constructions, interrogatives and some negation.
+
+English is the largest language by number of speakers,[14] and the third most-spoken native language in the world, after Standard Chinese and Spanish.[15] It is the most widely learned second language and is either the official language or one of the official languages in almost 60 sovereign states. There are more people who have learned it as a second language than there are native speakers. It is estimated that there are over 2 billion speakers of English.[16] English is the majority native language in the United States, the United Kingdom, Canada, Australia, New Zealand and Ireland, and it is widely spoken in some areas of the Caribbean, Africa and South Asia.[17] It is a co-official language of the United Nations, the European Union and many other world and regional international organisations. It is the most widely spoken Germanic language, accounting for at least 70% of speakers of this Indo-European branch. English speakers are called "Anglophones". Variability among the accents and dialects of English used in different countries and regions—in terms of phonetics and phonology, and sometimes also vocabulary, idioms, grammar, and spelling—does not typically prevent understanding by speakers of other dialects, although mutual unintelligibility can occur at extreme ends of the dialect continuum.
+
+English is an Indo-European language and belongs to the West Germanic group of the Germanic languages.[18] Old English originated from a Germanic tribal and linguistic continuum along the Frisian North Sea coast, whose languages gradually evolved into the Anglic languages in the British Isles, and into the Frisian languages and Low German/Low Saxon on the continent. The Frisian languages, which together with the Anglic languages form the Anglo-Frisian languages, are the closest living relatives of English. Low German/Low Saxon is also closely related, and sometimes English, the Frisian languages, and Low German are grouped together as the Ingvaeonic (North Sea Germanic) languages, though this grouping remains debated.[7] Old English evolved into Middle English, which in turn evolved into Modern English.[19] Particular dialects of Old and Middle English also developed into a number of other Anglic languages, including Scots[20] and the extinct Fingallian and Forth and Bargy (Yola) dialects of Ireland.[21]
+
+Like Icelandic and Faroese, the development of English in the British Isles isolated it from the continental Germanic languages and influences. It has since evolved considerably. English is not mutually intelligible with any continental Germanic language, differing in vocabulary, syntax, and phonology, although some of these, such as Dutch or Frisian, do show strong affinities with English, especially with its earlier stages.[22]
+
+Unlike Icelandic and Faroese, which were isolated, the development of English was influenced by a long series of invasions of the British Isles by other peoples and languages, particularly Old Norse and Norman French. These left a profound mark of their own on the language, so that English shows some similarities in vocabulary and grammar with many languages outside its linguistic clades—but it is not mutually intelligible with any of those languages either. Some scholars have argued that English can be considered a mixed language or a creole—a theory called the Middle English creole hypothesis. Although the great influence of these languages on the vocabulary and grammar of Modern English is widely acknowledged, most specialists in language contact do not consider English to be a true mixed language.[23][24]
+
+English is classified as a Germanic language because it shares innovations with other Germanic languages such as Dutch, German, and Swedish.[25] These shared innovations show that the languages have descended from a single common ancestor called Proto-Germanic. Some shared features of Germanic languages include the division of verbs into strong and weak classes, the use of modal verbs, and the sound changes affecting Proto-Indo-European consonants, known as Grimm's and Verner's laws. English is classified as an Anglo-Frisian language because Frisian and English share other features, such as the palatalisation of consonants that were velar consonants in Proto-Germanic (see Phonological history of Old English § Palatalization).[26]
+
+The earliest form of English is called Old English or Anglo-Saxon (c. 550–1066 CE). Old English developed from a set of West Germanic dialects, often grouped as Anglo-Frisian or North Sea Germanic, and originally spoken along the coasts of Frisia, Lower Saxony and southern Jutland by Germanic peoples known to the historical record as the Angles, Saxons, and Jutes.[27][28] From the 5th century CE, the Anglo-Saxons settled Britain as the Roman economy and administration collapsed. By the 7th century, the Germanic language of the Anglo-Saxons became dominant in Britain, replacing the languages of Roman Britain (43–409 CE): Common Brittonic, a Celtic language, and Latin, brought to Britain by the Roman occupation.[29][30][31] England and English (originally Ænglaland and Ænglisc) are named after the Angles.[32]
+
+Old English was divided into four dialects: the Anglian dialects (Mercian and Northumbrian) and the Saxon dialects, Kentish and West Saxon.[33] Through the educational reforms of King Alfred in the 9th century and the influence of the kingdom of Wessex, the West Saxon dialect became the standard written variety.[34] The epic poem Beowulf is written in West Saxon, and the earliest English poem, Cædmon's Hymn, is written in Northumbrian.[35] Modern English developed mainly from Mercian, but the Scots language developed from Northumbrian. A few short inscriptions from the early period of Old English were written using a runic script.[36] By the 6th century, a Latin alphabet was adopted, written with half-uncial letterforms. It included the runic letters wynn ⟨ƿ⟩ and thorn ⟨þ⟩, and the modified Latin letters eth ⟨ð⟩, and ash ⟨æ⟩.[36][37]
+
+Old English is essentially a distinct language from Modern English and is virtually impossible for 21st-century unstudied English-speakers to understand. Its grammar was similar to that of modern German, and its closest relative is Old Frisian. Nouns, adjectives, pronouns, and verbs had many more inflectional endings and forms, and word order was much freer than in Modern English. Modern English has case forms in pronouns (he, him, his) and has a few verb inflections (speak, speaks, speaking, spoke, spoken), but Old English had case endings in nouns as well, and verbs had more person and number endings.[38][39][40]
+
+The translation of Matthew 8:20 from 1000 CE shows examples of case endings (nominative plural, accusative plural, genitive singular) and a verb ending (present plural):
+
+John of Trevisa, ca. 1385[42]
+
+From the 8th to the 12th century, Old English gradually transformed through language contact into Middle English. Middle English is often arbitrarily defined as beginning with the conquest of England by William the Conqueror in 1066, but it developed further in the period from 1200–1450.
+
+First, the waves of Norse colonisation of northern parts of the British Isles in the 8th and 9th centuries put Old English into intense contact with Old Norse, a North Germanic language. Norse influence was strongest in the north-eastern varieties of Old English spoken in the Danelaw area around York, which was the centre of Norse colonisation; today these features are still particularly present in Scots and Northern English. However the centre of norsified English seems to have been in the Midlands around Lindsey, and after 920 CE when Lindsey was reincorporated into the Anglo-Saxon polity, Norse features spread from there into English varieties that had not been in direct contact with Norse speakers. An element of Norse influence that persists in all English varieties today is the group of pronouns beginning with th- (they, them, their) which replaced the Anglo-Saxon pronouns with h- (hie, him, hera).[43]
+
+With the Norman conquest of England in 1066, the now norsified Old English language was subject to contact with Old French, in particular with the Old Norman dialect. The Norman language in England eventually developed into Anglo-Norman.[9] Because Norman was spoken primarily by the elites and nobles, while the lower classes continued speaking Anglo-Saxon (English), the main influence of Norman was the introduction of a wide range of loanwords related to politics, legislation and prestigious social domains.[8] Middle English also greatly simplified the inflectional system, probably in order to reconcile Old Norse and Old English, which were inflectionally different but morphologically similar. The distinction between nominative and accusative cases was lost except in personal pronouns, the instrumental case was dropped, and the use of the genitive case was limited to indicating possession. The inflectional system regularised many irregular inflectional forms,[44] and gradually simplified the system of agreement, making word order less flexible.[45] In the Wycliffe Bible of the 1380s, the verse Matthew 8:20 was written:
+
+Here the plural suffix -n on the verb have is still retained, but none of the case endings on the nouns are present. By the 12th century Middle English was fully developed, integrating both Norse and French features; it continued to be spoken until the transition to early Modern English around 1500. Middle English literature includes Geoffrey Chaucer's The Canterbury Tales, and Malory's Le Morte d'Arthur. In the Middle English period, the use of regional dialects in writing proliferated, and dialect traits were even used for effect by authors such as Chaucer.[47]
+
+The next period in the history of English was Early Modern English (1500–1700). Early Modern English was characterised by the Great Vowel Shift (1350–1700), inflectional simplification, and linguistic standardisation.
+
+The Great Vowel Shift affected the stressed long vowels of Middle English. It was a chain shift, meaning that each shift triggered a subsequent shift in the vowel system. Mid and open vowels were raised, and close vowels were broken into diphthongs. For example, the word bite was originally pronounced as the word beet is today, and the second vowel in the word about was pronounced as the word boot is today. The Great Vowel Shift explains many irregularities in spelling since English retains many spellings from Middle English, and it also explains why English vowel letters have very different pronunciations from the same letters in other languages.[48][49]
+
+English began to rise in prestige, relative to Norman French, during the reign of Henry V. Around 1430, the Court of Chancery in Westminster began using English in its official documents, and a new standard form of Middle English, known as Chancery Standard, developed from the dialects of London and the East Midlands. In 1476, William Caxton introduced the printing press to England and began publishing the first printed books in London, expanding the influence of this form of English.[50] Literature from the Early Modern period includes the works of William Shakespeare and the translation of the Bible commissioned by King James I. Even after the vowel shift the language still sounded different from Modern English: for example, the consonant clusters /kn ɡn sw/ in knight, gnat, and sword were still pronounced. Many of the grammatical features that a modern reader of Shakespeare might find quaint or archaic represent the distinct characteristics of Early Modern English.[51]
+
+In the 1611 King James Version of the Bible, written in Early Modern English, Matthew 8:20 says:
+
+This exemplifies the loss of case and its effects on sentence structure (replacement with Subject-Verb-Object word order, and the use of of instead of the non-possessive genitive), and the introduction of loanwords from French (ayre) and word replacements (bird originally meaning "nestling" had replaced OE fugol).[52]
+
+By the late 18th century, the British Empire had spread English through its colonies and geopolitical dominance. Commerce, science and technology, diplomacy, art, and formal education all contributed to English becoming the first truly global language. English also facilitated worldwide international communication.[53][12] England continued to form new colonies, and these later developed their own norms for speech and writing. English was adopted in parts of North America, parts of Africa, Australasia, and many other regions. When they obtained political independence, some of the newly independent nations that had multiple indigenous languages opted to continue using English as the official language to avoid the political and other difficulties inherent in promoting any one indigenous language above the others.[54][55][56] In the 20th century the growing economic and cultural influence of the United States and its status as a superpower following the Second World War has, along with worldwide broadcasting in English by the BBC[57] and other broadcasters, caused the language to spread across the planet much faster.[58][59] In the 21st century, English is more widely spoken and written than any language has ever been.[60]
+
+As Modern English developed, explicit norms for standard usage were published, and spread through official media such as public education and state-sponsored publications. In 1755 Samuel Johnson published his A Dictionary of the English Language which introduced standard spellings of words and usage norms. In 1828, Noah Webster published the American Dictionary of the English language to try to establish a norm for speaking and writing American English that was independent of the British standard. Within Britain, non-standard or lower class dialect features were increasingly stigmatised, leading to the quick spread of the prestige varieties among the middle classes.[61]
+
+In modern English, the loss of grammatical case is almost complete (it is now only found in pronouns, such as he and him, she and her, who and whom), and SVO word-order is mostly fixed.[61] Some changes, such as the use of do-support have become universalised. (Earlier English did not use the word "do" as a general auxiliary as Modern English does; at first it was only used in question constructions, and even then was not obligatory.[62] Now, do-support with the verb have is becoming increasingly standardised.) The use of progressive forms in -ing, appears to be spreading to new constructions, and forms such as had been being built are becoming more common. Regularisation of irregular forms also slowly continues (e.g. dreamed instead of dreamt), and analytical alternatives to inflectional forms are becoming more common (e.g. more polite instead of politer). British English is also undergoing change under the influence of American English, fuelled by the strong presence of American English in the media and the prestige associated with the US as a world power.[63][64][65]
+
+
+
+As of 2016[update], 400 million people spoke English as their first language, and 1.1 billion spoke it as a secondary language.[66] English is the largest language by number of speakers. English is spoken by communities on every continent and on islands in all the major oceans.[67]
+
+The countries where English is spoken can be grouped into different categories according to how English is used in each country. The "inner circle"[68] countries with many native speakers of English share an international standard of written English and jointly influence speech norms for English around the world. English does not belong to just one country, and it does not belong solely to descendants of English settlers. English is an official language of countries populated by few descendants of native speakers of English. It has also become by far the most important language of international communication when people who share no native language meet anywhere in the world.
+
+Braj Kachru distinguishes countries where English is spoken with a three circles model.[68] In his model,
+
+Kachru bases his model on the history of how English spread in different countries, how users acquire English, and the range of uses English has in each country. The three circles change membership over time.[69]
+
+Countries with large communities of native speakers of English (the inner circle) include Britain, the United States, Australia, Canada, Ireland, and New Zealand, where the majority speaks English, and South Africa, where a significant minority speaks English. The countries with the most native English speakers are, in descending order, the United States (at least 231 million),[70] the United Kingdom (60 million),[71][72][73] Canada (19 million),[74] Australia (at least 17 million),[75] South Africa (4.8 million),[76] Ireland (4.2 million), and New Zealand (3.7 million).[77] In these countries, children of native speakers learn English from their parents, and local people who speak other languages and new immigrants learn English to communicate in their neighbourhoods and workplaces.[78] The inner-circle countries provide the base from which English spreads to other countries in the world.[69]
+
+Estimates of the numbers of second language and foreign-language English speakers vary greatly from 470 million to more than 1 billion, depending on how proficiency is defined.[17] Linguist David Crystal estimates that non-native speakers now outnumber native speakers by a ratio of 3 to 1.[79] In Kachru's three-circles model, the "outer circle" countries are countries such as the Philippines,[80] Jamaica,[81] India, Pakistan, Singapore,[82] Malaysia and Nigeria[83][84] with a much smaller proportion of native speakers of English but much use of English as a second language for education, government, or domestic business, and its routine use for school instruction and official interactions with the government.[85]
+
+Those countries have millions of native speakers of dialect continua ranging from an English-based creole to a more standard version of English. They have many more speakers of English who acquire English as they grow up through day-to-day use and listening to broadcasting, especially if they attend schools where English is the medium of instruction. Varieties of English learned by non-native speakers born to English-speaking parents may be influenced, especially in their grammar, by the other languages spoken by those learners.[78] Most of those varieties of English include words little used by native speakers of English in the inner-circle countries,[78] and they may show grammatical and phonological differences from inner-circle varieties as well. The standard English of the inner-circle countries is often taken as a norm for use of English in the outer-circle countries.[78]
+
+In the three-circles model, countries such as Poland, China, Brazil, Germany, Japan, Indonesia, Egypt, and other countries where English is taught as a foreign language, make up the "expanding circle".[86] The distinctions between English as a first language, as a second language, and as a foreign language are often debatable and may change in particular countries over time.[85] For example, in the Netherlands and some other countries of Europe, knowledge of English as a second language is nearly universal, with over 80 percent of the population able to use it,[87] and thus English is routinely used to communicate with foreigners and often in higher education. In these countries, although English is not used for government business, its widespread use puts them at the boundary between the "outer circle" and "expanding circle". English is unusual among world languages in how many of its users are not native speakers but speakers of English as a second or foreign language.[88]
+
+Many users of English in the expanding circle use it to communicate with other people from the expanding circle, so that interaction with native speakers of English plays no part in their decision to use English.[89] Non-native varieties of English are widely used for international communication, and speakers of one such variety often encounter features of other varieties.[90] Very often today a conversation in English anywhere in the world may include no native speakers of English at all, even while including speakers from several different countries.[91]
+
+Pie chart showing the percentage of native English speakers living in "inner circle" English-speaking countries. Native speakers are now substantially outnumbered worldwide by second-language speakers of English (not counted in this chart).
+
+English is a pluricentric language, which means that no one national authority sets the standard for use of the language.[92][93][94][95] But English is not a divided language,[96] despite a long-standing joke originally attributed to George Bernard Shaw that the United Kingdom and the United States are "two countries separated by a common language".[97] Spoken English, for example English used in broadcasting, generally follows national pronunciation standards that are also established by custom rather than by regulation. International broadcasters are usually identifiable as coming from one country rather than another through their accents,[98] but newsreader scripts are also composed largely in international standard written English. The norms of standard written English are maintained purely by the consensus of educated English-speakers around the world, without any oversight by any government or international organisation.[99]
+
+American listeners generally readily understand most British broadcasting, and British listeners readily understand most American broadcasting. Most English speakers around the world can understand radio programmes, television programmes, and films from many parts of the English-speaking world.[100] Both standard and non-standard varieties of English can include both formal or informal styles, distinguished by word choice and syntax and use both technical and non-technical registers.[101]
+
+The settlement history of the English-speaking inner circle countries outside Britain helped level dialect distinctions and produce koineised forms of English in South Africa, Australia, and New Zealand.[102] The majority of immigrants to the United States without British ancestry rapidly adopted English after arrival. Now the majority of the United States population are monolingual English speakers,[103][70] and English has been given official or co-official status by 30 of the 50 state governments, as well as all five territorial governments of the US, though there has never been an official language at the Federal level.[104][105]
+
+English has ceased to be an "English language" in the sense of belonging only to people who are ethnically English.[106][107] Use of English is growing country-by-country internally and for international communication. Most people learn English for practical rather than ideological reasons.[108] Many speakers of English in Africa have become part of an "Afro-Saxon" language community that unites Africans from different countries.[109]
+
+As decolonisation proceeded throughout the British Empire in the 1950s and 1960s, former colonies often did not reject English but rather continued to use it as independent countries setting their own language policies.[55][56][110] For example, the view of the English language among many Indians has gone from associating it with colonialism to associating it with economic progress, and English continues to be an official language of India.[111] English is also widely used in media and literature, and the number of English language books published annually in India is the third largest in the world after the US and UK.[112] However English is rarely spoken as a first language, numbering only around a couple hundred-thousand people, and less than 5% of the population speak fluent English in India.[113][114] David Crystal claimed in 2004 that, combining native and non-native speakers, India now has more people who speak or understand English than any other country in the world,[115] but the number of English speakers in India is very uncertain, with most scholars concluding that the United States still has more speakers of English than India.[116]
+
+Modern English, sometimes described as the first global lingua franca,[58][117] is also regarded as the first world language.[118][119] English is the world's most widely used language in newspaper publishing, book publishing, international telecommunications, scientific publishing, international trade, mass entertainment, and diplomacy.[119] English is, by international treaty, the basis for the required controlled natural languages[120] Seaspeak and Airspeak, used as international languages of seafaring[121] and aviation.[122] English used to have parity with French and German in scientific research, but now it dominates that field.[123] It achieved parity with French as a language of diplomacy at the Treaty of Versailles negotiations in 1919.[124] By the time of the foundation of the United Nations at the end of World War II, English had become pre-eminent[125] and is now the main worldwide language of diplomacy and international relations.[126] It is one of six official languages of the United Nations.[127] Many other worldwide international organisations, including the International Olympic Committee, specify English as a working language or official language of the organisation.
+
+Many regional international organisations such as the European Free Trade Association, Association of Southeast Asian Nations (ASEAN),[59] and Asia-Pacific Economic Cooperation (APEC) set English as their organisation's sole working language even though most members are not countries with a majority of native English speakers. While the European Union (EU) allows member states to designate any of the national languages as an official language of the Union, in practice English is the main working language of EU organisations.[128]
+
+Although in most countries English is not an official language, it is currently the language most often taught as a foreign language.[58][59] In the countries of the EU, English is the most widely spoken foreign language in nineteen of the twenty-five member states where it is not an official language (that is, the countries other than Ireland and Malta). In a 2012 official Eurobarometer poll (conducted when the UK was still a member of the EU), 38 percent of the EU respondents outside the countries where English is an official language said they could speak English well enough to have a conversation in that language. The next most commonly mentioned foreign language, French (which is the most widely known foreign language in the UK and Ireland), could be used in conversation by 12 percent of respondents.[129]
+
+A working knowledge of English has become a requirement in a number of occupations and professions such as medicine[130] and computing. English has become so important in scientific publishing that more than 80 percent of all scientific journal articles indexed by Chemical Abstracts in 1998 were written in English, as were 90 percent of all articles in natural science publications by 1996 and 82 percent of articles in humanities publications by 1995.[131]
+
+International communities such as international business people may use English as an auxiliary language, with an emphasis on vocabulary suitable for their domain of interest. This has led some scholars to develop the study of English as an auxiliary language. The trademarked Globish uses a relatively small subset of English vocabulary (about 1500 words, designed to represent the highest use in international business English) in combination with the standard English grammar.[132] Other examples include Simple English.
+
+The increased use of the English language globally has had an effect on other languages, leading to some English words being assimilated into the vocabularies of other languages. This influence of English has led to concerns about language death,[133] and to claims of linguistic imperialism,[134] and has provoked resistance to the spread of English; however the number of speakers continues to increase because many people around the world think that English provides them with opportunities for better employment and improved lives.[135]
+
+Although some scholars[who?] mention a possibility of future divergence of English dialects into mutually unintelligible languages, most think a more likely outcome is that English will continue to function as a koineised language in which the standard form unifies speakers from around the world.[136] English is used as the language for wider communication in countries around the world.[137] Thus English has grown in worldwide use much more than any constructed language proposed as an international auxiliary language, including Esperanto.[138][139]
+
+The phonetics and phonology of the English language differ from one dialect to another, usually without interfering with mutual communication. Phonological variation affects the inventory of phonemes (i.e. speech sounds that distinguish meaning), and phonetic variation consists in differences in pronunciation of the phonemes. [140] This overview mainly describes the standard pronunciations of the United Kingdom and the United States: Received Pronunciation (RP) and General American (GA). (See § Dialects, accents, and varieties, below.)
+
+The phonetic symbols used below are from the International Phonetic Alphabet (IPA).[141][142][143]
+
+Most English dialects share the same 24 consonant phonemes. The consonant inventory shown below is valid for California English,[144] and for RP.[145]
+
+* Conventionally transcribed /r/
+
+In the table, when obstruents (stops, affricates, and fricatives) appear in pairs, such as /p b/, /tʃ dʒ/, and /s z/, the first is fortis (strong) and the second is lenis (weak). Fortis obstruents, such as /p tʃ s/ are pronounced with more muscular tension and breath force than lenis consonants, such as /b dʒ z/, and are always voiceless. Lenis consonants are partly voiced at the beginning and end of utterances, and fully voiced between vowels. Fortis stops such as /p/ have additional articulatory or acoustic features in most dialects: they are aspirated [pʰ] when they occur alone at the beginning of a stressed syllable, often unaspirated in other cases, and often unreleased [p̚] or pre-glottalised [ʔp] at the end of a syllable. In a single-syllable word, a vowel before a fortis stop is shortened: thus nip has a noticeably shorter vowel (phonetically, but not phonemically) than nib [nɪˑb̥] (see below).[146]
+
+In RP, the lateral approximant /l/, has two main allophones (pronunciation variants): the clear or plain [l], as in light, and the dark or velarised [ɫ], as in full.[147] GA has dark l in most cases.[148]
+
+All sonorants (liquids /l, r/ and nasals /m, n, ŋ/) devoice when following a voiceless obstruent, and they are syllabic when following a consonant at the end of a word.[149]
+
+The pronunciation of vowels varies a great deal between dialects and is one of the most detectable aspects of a speaker's accent. The table below lists the vowel phonemes in Received Pronunciation (RP) and General American (GA), with examples of words in which they occur from lexical sets compiled by linguists. The vowels are represented with symbols from the International Phonetic Alphabet; those given for RP are standard in British dictionaries and other publications.[150]
+
+In RP, vowel length is phonemic; long vowels are marked with a triangular colon ⟨ː⟩ in the table above, such as the vowel of need [niːd] as opposed to bid [bɪd]. In GA, vowel length is non-distinctive.
+
+In both RP and GA, vowels are phonetically shortened before fortis consonants in the same syllable, like /t tʃ f/, but not before lenis consonants like /d dʒ v/ or in open syllables: thus, the vowels of rich [rɪtʃ], neat [nit], and safe [seɪ̯f] are noticeably shorter than the vowels of ridge [rɪˑdʒ], need [niˑd], and save [seˑɪ̯v], and the vowel of light [laɪ̯t] is shorter than that of lie [laˑɪ̯]. Because lenis consonants are frequently voiceless at the end of a syllable, vowel length is an important cue as to whether the following consonant is lenis or fortis.[151]
+
+The vowel /ə/ only occurs in unstressed syllables and is more open in quality in stem-final positions.[152][153] Some dialects do not contrast /ɪ/ and /ə/ in unstressed positions, so that rabbit and abbot rhyme and Lenin and Lennon are homophonous, a dialect feature called weak vowel merger.[154] GA /ɜr/ and /ər/ are realised as an r-coloured vowel [ɚ], as in further [ˈfɚðɚ] (phonemically /ˈfɜrðər/), which in RP is realised as [ˈfəːðə] (phonemically /ˈfɜːðə/).[155]
+
+An English syllable includes a syllable nucleus consisting of a vowel sound. Syllable onset and coda (start and end) are optional. A syllable can start with up to three consonant sounds, as in sprint /sprɪnt/, and end with up to four, as in texts /teksts/. This gives an English syllable the following structure, (CCC)V(CCCC) where C represents a consonant and V a vowel; the word strengths /strɛŋkθs/ is thus an example of the most complex syllable possible in English. The consonants that may appear together in onsets or codas are restricted, as is the order in which they may appear. Onsets can only have four types of consonant clusters: a stop and approximant, as in play; a voiceless fricative and approximant, as in fly or sly; s and a voiceless stop, as in stay; and s, a voiceless stop, and an approximant, as in string.[156] Clusters of nasal and stop are only allowed in codas. Clusters of obstruents always agree in voicing, and clusters of sibilants and of plosives with the same point of articulation are prohibited. Furthermore, several consonants have limited distributions: /h/ can only occur in syllable-initial position, and /ŋ/ only in syllable-final position.[157]
+
+Stress plays an important role in English. Certain syllables are stressed, while others are unstressed. Stress is a combination of duration, intensity, vowel quality, and sometimes changes in pitch. Stressed syllables are pronounced longer and louder than unstressed syllables, and vowels in unstressed syllables are frequently reduced while vowels in stressed syllables are not.[158] Some words, primarily short function words but also some modal verbs such as can, have weak and strong forms depending on whether they occur in stressed or non-stressed position within a sentence.
+
+Stress in English is phonemic, and some pairs of words are distinguished by stress. For instance, the word contract is stressed on the first syllable (/ˈkɒntrækt/ KON-trakt) when used as a noun, but on the last syllable (/kənˈtrækt/ kən-TRAKT) for most meanings (for example, "reduce in size") when used as a verb.[159][160][161] Here stress is connected to vowel reduction: in the noun "contract" the first syllable is stressed and has the unreduced vowel /ɒ/, but in the verb "contract" the first syllable is unstressed and its vowel is reduced to /ə/. Stress is also used to distinguish between words and phrases, so that a compound word receives a single stress unit, but the corresponding phrase has two: e.g. a burnout (/ˈbɜːrnaʊt/) versus to burn out (/ˈbɜːrn ˈaʊt/), and a hotdog (/ˈhɒtdɒɡ/) versus a hot dog (/ˈhɒt ˈdɒɡ/).[162]
+
+In terms of rhythm, English is generally described as a stress-timed language, meaning that the amount of time between stressed syllables tends to be equal.[163] Stressed syllables are pronounced longer, but unstressed syllables (syllables between stresses) are shortened. Vowels in unstressed syllables are shortened as well, and vowel shortening causes changes in vowel quality: vowel reduction.[164]
+
+Varieties of English vary the most in pronunciation of vowels. The best known national varieties used as standards for education in non-English-speaking countries are British (BrE) and American (AmE). Countries such as Canada, Australia, Ireland, New Zealand and South Africa have their own standard varieties which are less often used as standards for education internationally. Some differences between the various dialects are shown in the table "Varieties of Standard English and their features".[165]
+
+English has undergone many historical sound changes, some of them affecting all varieties, and others affecting only a few. Most standard varieties are affected by the Great Vowel Shift, which changed the pronunciation of long vowels, but a few dialects have slightly different results. In North America, a number of chain shifts such as the Northern Cities Vowel Shift and Canadian Shift have produced very different vowel landscapes in some regional accents.[166][167]
+
+Some dialects have fewer or more consonant phonemes and phones than the standard varieties. Some conservative varieties like Scottish English have a voiceless [ʍ] sound in whine that contrasts with the voiced [w] in wine, but most other dialects pronounce both words with voiced [w], a dialect feature called wine–whine merger. The unvoiced velar fricative sound /x/ is found in Scottish English, which distinguishes loch /lɔx/ from lock /lɔk/. Accents like Cockney with "h-dropping" lack the glottal fricative /h/, and dialects with th-stopping and th-fronting like African American Vernacular and Estuary English do not have the dental fricatives /θ, ð/, but replace them with dental or alveolar stops /t, d/ or labiodental fricatives /f, v/.[168][169] Other changes affecting the phonology of local varieties are processes such as yod-dropping, yod-coalescence, and reduction of consonant clusters.[170]
+
+General American and Received Pronunciation vary in their pronunciation of historical /r/ after a vowel at the end of a syllable (in the syllable coda). GA is a rhotic dialect, meaning that it pronounces /r/ at the end of a syllable, but RP is non-rhotic, meaning that it loses /r/ in that position. English dialects are classified as rhotic or non-rhotic depending on whether they elide /r/ like RP or keep it like GA.[171]
+
+There is complex dialectal variation in words with the open front and open back vowels /æ ɑː ɒ ɔː/. These four vowels are only distinguished in RP, Australia, New Zealand and South Africa. In GA, these vowels merge to three /æ ɑ ɔ/,[172] and in Canadian English, they merge to two /æ ɑ/.[173] In addition, the words that have each vowel vary by dialect. The table "Dialects and open vowels" shows this variation with lexical sets in which these sounds occur.
+
+As is typical of an Indo-European language, English follows accusative morphosyntactic alignment. Unlike other Indo-European languages though, English has largely abandoned the inflectional case system in favor of analytic constructions. Only the personal pronouns retain morphological case more strongly than any other word class. English distinguishes at least seven major word classes: verbs, nouns, adjectives, adverbs, determiners (including articles), prepositions, and conjunctions. Some analyses add pronouns as a class separate from nouns, and subdivide conjunctions into subordinators and coordinators, and add the class of interjections.[174] English also has a rich set of auxiliary verbs, such as have and do, expressing the categories of mood and aspect. Questions are marked by do-support, wh-movement (fronting of question words beginning with wh-) and word order inversion with some verbs.[175]
+
+Some traits typical of Germanic languages persist in English, such as the distinction between irregularly inflected strong stems inflected through ablaut (i.e. changing the vowel of the stem, as in the pairs speak/spoke and foot/feet) and weak stems inflected through affixation (such as love/loved, hand/hands).[176] Vestiges of the case and gender system are found in the pronoun system (he/him, who/whom) and in the inflection of the copula verb to be.[176]
+
+The seven-word classes are exemplified in this sample sentence:[177]
+
+English nouns are only inflected for number and possession. New nouns can be formed through derivation or compounding. They are semantically divided into proper nouns (names) and common nouns. Common nouns are in turn divided into concrete and abstract nouns, and grammatically into count nouns and mass nouns.[178]
+
+Most count nouns are inflected for plural number through the use of the plural suffix -s, but a few nouns have irregular plural forms. Mass nouns can only be pluralised through the use of a count noun classifier, e.g. one loaf of bread, two loaves of bread.[179]
+
+Regular plural formation:
+
+Irregular plural formation:
+
+Possession can be expressed either by the possessive enclitic -s (also traditionally called a genitive suffix), or by the preposition of. Historically the -s possessive has been used for animate nouns, whereas the of possessive has been reserved for inanimate nouns. Today this distinction is less clear, and many speakers use -s also with inanimates. Orthographically the possessive -s is separated from the noun root with an apostrophe.[175]
+
+Possessive constructions:
+
+Nouns can form noun phrases (NPs) where they are the syntactic head of the words that depend on them such as determiners, quantifiers, conjunctions or adjectives.[180] Noun phrases can be short, such as the man, composed only of a determiner and a noun. They can also include modifiers such as adjectives (e.g. red, tall, all) and specifiers such as determiners (e.g. the, that). But they can also tie together several nouns into a single long NP, using conjunctions such as and, or prepositions such as with, e.g. the tall man with the long red trousers and his skinny wife with the spectacles (this NP uses conjunctions, prepositions, specifiers, and modifiers). Regardless of length, an NP functions as a syntactic unit.[175] For example, the possessive enclitic can, in cases which do not lead to ambiguity, follow the entire noun phrase, as in The President of India's wife, where the enclitic follows India and not President.
+
+The class of determiners is used to specify the noun they precede in terms of definiteness, where the marks a definite noun and a or an an indefinite one. A definite noun is assumed by the speaker to be already known by the interlocutor, whereas an indefinite noun is not specified as being previously known. Quantifiers, which include one, many, some and all, are used to specify the noun in terms of quantity or number. The noun must agree with the number of the determiner, e.g. one man (sg.) but all men (pl.). Determiners are the first constituents in a noun phrase.[181]
+
+Adjectives modify a noun by providing additional information about their referents. In English, adjectives come before the nouns they modify and after determiners.[182] In Modern English, adjectives are not inflected, and they do not agree in form with the noun they modify, as adjectives in most other Indo-European languages do. For example, in the phrases the slender boy, and many slender girls, the adjective slender does not change form to agree with either the number or gender of the noun.
+
+Some adjectives are inflected for degree of comparison, with the positive degree unmarked, the suffix -er marking the comparative, and -est marking the superlative: a small boy, the boy is smaller than the girl, that boy is the smallest. Some adjectives have irregular comparative and superlative forms, such as good, better, and best. Other adjectives have comparatives formed by periphrastic constructions, with the adverb more marking the comparative, and most marking the superlative: happier or more happy, the happiest or most happy.[183] There is some variation among speakers regarding which adjectives use inflected or periphrastic comparison, and some studies have shown a tendency for the periphrastic forms to become more common at the expense of the inflected form.[184]
+
+English pronouns conserve many traits of case and gender inflection. The personal pronouns retain a difference between subjective and objective case in most persons (I/me, he/him, she/her, we/us, they/them) as well as a gender and animateness distinction in the third person singular (distinguishing he/she/it). The subjective case corresponds to the Old English nominative case, and the objective case is used both in the sense of the previous accusative case (in the role of patient, or direct object of a transitive verb), and in the sense of the Old English dative case (in the role of a recipient or indirect object of a transitive verb).[185][186] Subjective case is used when the pronoun is the subject of a finite clause, and otherwise, the objective case is used.[187] While grammarians such as Henry Sweet[188] and Otto Jespersen[189] noted that the English cases did not correspond to the traditional Latin based system, some contemporary grammars, for example Huddleston & Pullum (2002), retain traditional labels for the cases, calling them nominative and accusative cases respectively.
+
+Possessive pronouns exist in dependent and independent forms; the dependent form functions as a determiner specifying a noun (as in my chair), while the independent form can stand alone as if it were a noun (e.g. the chair is mine).[190] The English system of grammatical person no longer has a distinction between formal and informal pronouns of address (the old 2nd person singular familiar pronoun thou acquired a pejorative or inferior tinge of meaning and was abandoned), and the forms for 2nd person plural and singular are identical except in the reflexive form. Some dialects have introduced innovative 2nd person plural pronouns such as y'all found in Southern American English and African American (Vernacular) English or youse found in Australian English and ye in Irish English.
+
+Pronouns are used to refer to entities deictically or anaphorically. A deictic pronoun points to some person or object by identifying it relative to the speech situation—for example, the pronoun I identifies the speaker, and the pronoun you, the addressee. Anaphoric pronouns such as that refer back to an entity already mentioned or assumed by the speaker to be known by the audience, for example in the sentence I already told you that. The reflexive pronouns are used when the oblique argument is identical to the subject of a phrase (e.g. "he sent it to himself" or "she braced herself for impact").[191]
+
+Prepositional phrases (PP) are phrases composed of a preposition and one or more nouns, e.g. with the dog, for my friend, to school, in England.[192] Prepositions have a wide range of uses in English. They are used to describe movement, place, and other relations between different entities, but they also have many syntactic uses such as introducing complement clauses and oblique arguments of verbs.[192] For example, in the phrase I gave it to him, the preposition to marks the recipient, or Indirect Object of the verb to give. Traditionally words were only considered prepositions if they governed the case of the noun they preceded, for example causing the pronouns to use the objective rather than subjective form, "with her", "to me", "for us". But some contemporary grammars such as that of Huddleston & Pullum (2002:598–600) no longer consider government of case to be the defining feature of the class of prepositions, rather defining prepositions as words that can function as the heads of prepositional phrases.
+
+English verbs are inflected for tense and aspect and marked for agreement with present-tense third-person singular subject. Only the copula verb to be is still inflected for agreement with the plural and first and second person subjects.[183] Auxiliary verbs such as have and be are paired with verbs in the infinitive, past, or progressive forms. They form complex tenses, aspects, and moods. Auxiliary verbs differ from other verbs in that they can be followed by the negation, and in that they can occur as the first constituent in a question sentence.[193][194]
+
+Most verbs have six inflectional forms. The primary forms are a plain present, a third-person singular present, and a preterite (past) form. The secondary forms are a plain form used for the infinitive, a gerund-participle and a past participle.[195] The copula verb to be is the only verb to retain some of its original conjugation, and takes different inflectional forms depending on the subject. The first-person present-tense form is am, the third person singular form is is, and the form are is used in the second-person singular and all three plurals. The only verb past participle is been and its gerund-participle is being.
+
+English has two primary tenses, past (preterit) and non-past. The preterit is inflected by using the preterit form of the verb, which for the regular verbs includes the suffix -ed, and for the strong verbs either the suffix -t or a change in the stem vowel. The non-past form is unmarked except in the third person singular, which takes the suffix -s.[193]
+
+English does not have a morphologised future tense.[196] Futurity of action is expressed periphrastically with one of the auxiliary verbs will or shall.[197] Many varieties also use a near future constructed with the phrasal verb be going to.[198]
+
+Further aspectual distinctions are encoded by the use of auxiliary verbs, primarily have and be, which encode the contrast between a perfect and non-perfect past tense (I have run vs. I was running), and compound tenses such as preterite perfect (I had been running) and present perfect (I have been running).[199]
+
+For the expression of mood, English uses a number of modal auxiliaries, such as can, may, will, shall and the past tense forms could, might, would, should. There is also a subjunctive and an imperative mood, both based on the plain form of the verb (i.e. without the third person singular -s), and which is used in subordinate clauses (e.g. subjunctive: It is important that he run every day; imperative Run!).[197]
+
+An infinitive form, that uses the plain form of the verb and the preposition to, is used for verbal clauses that are syntactically subordinate to a finite verbal clause. Finite verbal clauses are those that are formed around a verb in the present or preterit form. In clauses with auxiliary verbs, they are the finite verbs and the main verb is treated as a subordinate clause.[200] For example, he has to go where only the auxiliary verb have is inflected for time and the main verb to go is in the infinitive, or in a complement clause such as I saw him leave, where the main verb is to see which is in a preterite form, and leave is in the infinitive.
+
+English also makes frequent use of constructions traditionally called phrasal verbs, verb phrases that are made up of a verb root and a preposition or particle which follows the verb. The phrase then functions as a single predicate. In terms of intonation the preposition is fused to the verb, but in writing it is written as a separate word. Examples of phrasal verbs are to get up, to ask out, to back up, to give up, to get together, to hang out, to put up with, etc. The phrasal verb frequently has a highly idiomatic meaning that is more specialised and restricted than what can be simply extrapolated from the combination of verb and preposition complement (e.g. lay off meaning terminate someone's employment).[201] In spite of the idiomatic meaning, some grammarians, including Huddleston & Pullum (2002:274), do not consider this type of construction to form a syntactic constituent and hence refrain from using the term "phrasal verb". Instead, they consider the construction simply to be a verb with a prepositional phrase as its syntactic complement, i.e. he woke up in the morning and he ran up in the mountains are syntactically equivalent.
+
+The function of adverbs is to modify the action or event described by the verb by providing additional information about the manner in which it occurs.[175] Many adverbs are derived from adjectives by appending the suffix -ly. For example, in the phrase the woman walked quickly, the adverb quickly is derived in this way from the adjective quick. Some commonly used adjectives have irregular adverbial forms, such as good which has the adverbial form well.
+
+Modern English syntax language is moderately analytic.[202] It has developed features such as modal verbs and word order as resources for conveying meaning. Auxiliary verbs mark constructions such as questions, negative polarity, the passive voice and progressive aspect.
+
+English word order has moved from the Germanic verb-second (V2) word order to being almost exclusively subject–verb–object (SVO).[203] The combination of SVO order and use of auxiliary verbs often creates clusters of two or more verbs at the centre of the sentence, such as he had hoped to try to open it.
+
+In most sentences, English only marks grammatical relations through word order.[204] The subject constituent precedes the verb and the object constituent follows it. The example below demonstrates how the grammatical roles of each constituent is marked only by the position relative to the verb:
+
+An exception is found in sentences where one of the constituents is a pronoun, in which case it is doubly marked, both by word order and by case inflection, where the subject pronoun precedes the verb and takes the subjective case form, and the object pronoun follows the verb and takes the objective case form.[205] The example below demonstrates this double marking in a sentence where both object and subject is represented with a third person singular masculine pronoun:
+
+Indirect objects (IO) of ditransitive verbs can be placed either as the first object in a double object construction (S V IO O), such as I gave Jane the book or in a prepositional phrase, such as I gave the book to Jane.[206]
+
+In English a sentence may be composed of one or more clauses, that may, in turn, be composed of one or more phrases (e.g. Noun Phrases, Verb Phrases, and Prepositional Phrases). A clause is built around a verb and includes its constituents, such as any NPs and PPs. Within a sentence, there is always at least one main clause (or matrix clause) whereas other clauses are subordinate to a main clause. Subordinate clauses may function as arguments of the verb in the main clause. For example, in the phrase I think (that) you are lying, the main clause is headed by the verb think, the subject is I, but the object of the phrase is the subordinate clause (that) you are lying. The subordinating conjunction that shows that the clause that follows is a subordinate clause, but it is often omitted.[207] Relative clauses are clauses that function as a modifier or specifier to some constituent in the main clause: For example, in the sentence I saw the letter that you received today, the relative clause that you received today specifies the meaning of the word letter, the object of the main clause. Relative clauses can be introduced by the pronouns who, whose, whom and which as well as by that (which can also be omitted.)[208] In contrast to many other Germanic languages there is no major differences between word order in main and subordinate clauses.[209]
+
+English syntax relies on auxiliary verbs for many functions including the expression of tense, aspect, and mood. Auxiliary verbs form main clauses, and the main verbs function as heads of a subordinate clause of the auxiliary verb. For example, in the sentence the dog did not find its bone, the clause find its bone is the complement of the negated verb did not. Subject–auxiliary inversion is used in many constructions, including focus, negation, and interrogative constructions.
+
+The verb do can be used as an auxiliary even in simple declarative sentences, where it usually serves to add emphasis, as in "I did shut the fridge." However, in the negated and inverted clauses referred to above, it is used because the rules of English syntax permit these constructions only when an auxiliary is present. Modern English does not allow the addition of the negating adverb not to an ordinary finite lexical verb, as in *I know not—it can only be added to an auxiliary (or copular) verb, hence if there is no other auxiliary present when negation is required, the auxiliary do is used, to produce a form like I do not (don't) know. The same applies in clauses requiring inversion, including most questions—inversion must involve the subject and an auxiliary verb, so it is not possible to say *Know you him?; grammatical rules require Do you know him?[210]
+
+Negation is done with the adverb not, which precedes the main verb and follows an auxiliary verb. A contracted form of not -n't can be used as an enclitic attaching to auxiliary verbs and to the copula verb to be. Just as with questions, many negative constructions require the negation to occur with do-support, thus in Modern English I don't know him is the correct answer to the question Do you know him?, but not *I know him not, although this construction may be found in older English.[211]
+
+Passive constructions also use auxiliary verbs. A passive construction rephrases an active construction in such a way that the object of the active phrase becomes the subject of the passive phrase, and the subject of the active phrase is either omitted or demoted to a role as an oblique argument introduced in a prepositional phrase. They are formed by using the past participle either with the auxiliary verb to be or to get, although not all varieties of English allow the use of passives with get. For example, putting the sentence she sees him into the passive becomes he is seen (by her), or he gets seen (by her).[212]
+
+Both yes–no questions and wh-questions in English are mostly formed using subject–auxiliary inversion (Am I going tomorrow?, Where can we eat?), which may require do-support (Do you like her?, Where did he go?). In most cases, interrogative words (wh-words; e.g. what, who, where, when, why, how) appear in a fronted position. For example, in the question What did you see?, the word what appears as the first constituent despite being the grammatical object of the sentence. (When the wh-word is the subject or forms part of the subject, no inversion occurs: Who saw the cat?.) Prepositional phrases can also be fronted when they are the question's theme, e.g. To whose house did you go last night?. The personal interrogative pronoun who is the only interrogative pronoun to still show inflection for case, with the variant whom serving as the objective case form, although this form may be going out of use in many contexts.[213]
+
+While English is a subject-prominent language, at the discourse level it tends to use a topic-comment structure, where the known information (topic) precedes the new information (comment). Because of the strict SVO syntax, the topic of a sentence generally has to be the grammatical subject of the sentence. In cases where the topic is not the grammatical subject of the sentence, frequently the topic is promoted to subject position through syntactic means. One way of doing this is through a passive construction, the girl was stung by the bee. Another way is through a cleft sentence where the main clause is demoted to be a complement clause of a copula sentence with a dummy subject such as it or there, e.g. it was the girl that the bee stung, there was a girl who was stung by a bee.[214] Dummy subjects are also used in constructions where there is no grammatical subject such as with impersonal verbs (e.g., it is raining) or in existential clauses (there are many cars on the street). Through the use of these complex sentence constructions with informationally vacuous subjects, English is able to maintain both a topic-comment sentence structure and a SVO syntax.
+
+Focus constructions emphasise a particular piece of new or salient information within a sentence, generally through allocating the main sentence level stress on the focal constituent. For example, the girl was stung by a bee (emphasising it was a bee and not, for example, a wasp that stung her), or The girl was stung by a bee (contrasting with another possibility, for example that it was the boy).[215] Topic and focus can also be established through syntactic dislocation, either preposing or postposing the item to be focused on relative to the main clause. For example, That girl over there, she was stung by a bee, emphasises the girl by preposition, but a similar effect could be achieved by postposition, she was stung by a bee, that girl over there, where reference to the girl is established as an "afterthought".[216]
+
+Cohesion between sentences is achieved through the use of deictic pronouns as anaphora (e.g. that is exactly what I mean where that refers to some fact known to both interlocutors, or then used to locate the time of a narrated event relative to the time of a previously narrated event).[217] Discourse markers such as oh, so or well, also signal the progression of ideas between sentences and help to create cohesion. Discourse markers are often the first constituents in sentences. Discourse markers are also used for stance taking in which speakers position themselves in a specific attitude towards what is being said, for example, no way is that true! (the idiomatic marker no way! expressing disbelief), or boy! I'm hungry (the marker boy expressing emphasis). While discourse markers are particularly characteristic of informal and spoken registers of English, they are also used in written and formal registers.[218]
+
+English is a rich language in terms of vocabulary, containing more synonyms than any other language.[134] There are words which appear on the surface to mean exactly the same thing but which, in fact, have slightly different shades of meaning and must be chosen appropriately if a speaker wants to convey precisely the message intended. It is generally stated that English has around 170,000 words, or 220,000 if obsolete words are counted; this estimate is based on the last full edition of the Oxford English Dictionary from 1989.[219] Over half of these words are nouns, a quarter adjectives, and a seventh verbs. There is one count that puts the English vocabulary at about 1 million words—but that count presumably includes words such as Latin species names, scientific terminology, botanical terms, prefixed and suffixed words, jargon, foreign words of extremely limited English use, and technical acronyms.[220]
+
+Due to its status as an international language, English adopts foreign words quickly, and borrows vocabulary from many other sources. Early studies of English vocabulary by lexicographers, the scholars who formally study vocabulary, compile dictionaries, or both, were impeded by a lack of comprehensive data on actual vocabulary in use from good-quality linguistic corpora,[221] collections of actual written texts and spoken passages. Many statements published before the end of the 20th century about the growth of English vocabulary over time, the dates of first use of various words in English, and the sources of English vocabulary will have to be corrected as new computerised analysis of linguistic corpus data becomes available.[220][222]
+
+English forms new words from existing words or roots in its vocabulary through a variety of processes. One of the most productive processes in English is conversion,[223] using a word with a different grammatical role, for example using a noun as a verb or a verb as a noun. Another productive word-formation process is nominal compounding,[220][222] producing compound words such as babysitter or ice cream or homesick.[223] A process more common in Old English than in Modern English, but still productive in Modern English, is the use of derivational suffixes (-hood, -ness, -ing, -ility) to derive new words from existing words (especially those of Germanic origin) or stems (especially for words of Latin or Greek origin).
+
+Formation of new words, called neologisms, based on Greek and/or Latin roots (for example television or optometry) is a highly productive process in English and in most modern European languages, so much so that it is often difficult to determine in which language a neologism originated. For this reason, lexicographer Philip Gove attributed many such words to the "international scientific vocabulary" (ISV) when compiling Webster's Third New International Dictionary (1961). Another active word-formation process in English is acronyms,[224] words formed by pronouncing as a single word abbreviations of longer phrases (e.g. NATO, laser).
+
+Source languages of English vocabulary[6][225]
+
+English, besides forming new words from existing words and their roots, also borrows words from other languages. This adoption of words from other languages is commonplace in many world languages, but English has been especially open to borrowing of foreign words throughout the last 1,000 years.[226] The most commonly used words in English are West Germanic.[227] The words in English learned first by children as they learn to speak, particularly the grammatical words that dominate the word count of both spoken and written texts, are mainly the Germanic words inherited from the earliest periods of the development of Old English.[220]
+
+But one of the consequences of long language contact between French and English in all stages of their development is that the vocabulary of English has a very high percentage of "Latinate" words (derived from French, especially, and also from other Romance languages and Latin). French words from various periods of the development of French now make up one-third of the vocabulary of English.[228] Linguist Anthony Lacoudre estimated that over 40,000 English words are of French origin and may be understood without orthographical change by French speakers.[229] Words of Old Norse origin have entered the English language primarily from the contact between Old Norse and Old English during colonisation of eastern and northern England. Many of these words are part of English core vocabulary, such as egg and knife.[230]
+
+English has also borrowed many words directly from Latin, the ancestor of the Romance languages, during all stages of its development.[222][220] Many of these words had earlier been borrowed into Latin from Greek. Latin or Greek are still highly productive sources of stems used to form vocabulary of subjects learned in higher education such as the sciences, philosophy, and mathematics.[231] English continues to gain new loanwords and calques ("loan translations") from languages all over the world, and words from languages other than the ancestral Anglo-Saxon language make up about 60% of the vocabulary of English.[232]
+
+English has formal and informal speech registers; informal registers, including child-directed speech, tend to be made up predominantly of words of Anglo-Saxon origin, while the percentage of vocabulary that is of Latinate origin is higher in legal, scientific, and academic texts.[233][234]
+
+English has had a strong influence on the vocabulary of other languages.[228][235] The influence of English comes from such factors as opinion leaders in other countries knowing the English language, the role of English as a world lingua franca, and the large number of books and films that are translated from English into other languages.[236] That pervasive use of English leads to a conclusion in many places that English is an especially suitable language for expressing new ideas or describing new technologies. Among varieties of English, it is especially American English that influences other languages.[237] Some languages, such as Chinese, write words borrowed from English mostly as calques, while others, such as Japanese, readily take in English loanwords written in sound-indicating script.[238] Dubbed films and television programmes are an especially fruitful source of English influence on languages in Europe.[238]
+
+Since the ninth century, English has been written in a Latin alphabet (also called Roman alphabet). Earlier Old English texts in Anglo-Saxon runes are only short inscriptions. The great majority of literary works in Old English that survive to today are written in the Roman alphabet.[36] The modern English alphabet contains 26 letters of the Latin script: a, b, c, d, e, f, g, h, i, j, k, l, m, n, o, p, q, r, s, t, u, v, w, x, y, z (which also have capital forms: A, B, C, D, E, F, G, H, I, J, K, L, M, N, O, P, Q, R, S, T, U, V, W, X, Y, Z).
+
+The spelling system, or orthography, of English is multi-layered, with elements of French, Latin, and Greek spelling on top of the native Germanic system.[239] Further complications have arisen through sound changes with which the orthography has not kept pace.[48] Compared to European languages for which official organisations have promoted spelling reforms, English has spelling that is a less consistent indicator of pronunciation, and standard spellings of words that are more difficult to guess from knowing how a word is pronounced.[240] There are also systematic spelling differences between British and American English. These situations have prompted proposals for spelling reform in English.[241]
+
+Although letters and speech sounds do not have a one-to-one correspondence in standard English spelling, spelling rules that take into account syllable structure, phonetic changes in derived words, and word accent are reliable for most English words.[242] Moreover, standard English spelling shows etymological relationships between related words that would be obscured by a closer correspondence between pronunciation and spelling, for example the words photograph, photography, and photographic,[242] or the words electricity and electrical. While few scholars agree with Chomsky and Halle (1968) that conventional English orthography is "near-optimal",[239] there is a rationale for current English spelling patterns.[243] The standard orthography of English is the most widely used writing system in the world.[244] Standard English spelling is based on a graphomorphemic segmentation of words into written clues of what meaningful units make up each word.[245]
+
+Readers of English can generally rely on the correspondence between spelling and pronunciation to be fairly regular for letters or digraphs used to spell consonant sounds. The letters b, d, f, h, j, k, l, m, n, p, r, s, t, v, w, y, z represent, respectively, the phonemes /b, d, f, h, dʒ, k, l, m, n, p, r, s, t, v, w, j, z/. The letters c and g normally represent /k/ and /ɡ/, but there is also a soft c pronounced /s/, and a soft g pronounced /dʒ/. The differences in the pronunciations of the letters c and g are often signalled by the following letters in standard English spelling. Digraphs used to represent phonemes and phoneme sequences include ch for /tʃ/, sh for /ʃ/, th for /θ/ or /ð/, ng for /ŋ/, qu for /kw/, and ph for /f/ in Greek-derived words. The single letter x is generally pronounced as /z/ in word-initial position and as /ks/ otherwise. There are exceptions to these generalisations, often the result of loanwords being spelled according to the spelling patterns of their languages of origin[242] or residues of proposals by scholars in the early period of Modern English to follow the spelling patterns of Latin for English words of Germanic origin.[246]
+
+For the vowel sounds of the English language, however, correspondences between spelling and pronunciation are more irregular. There are many more vowel phonemes in English than there are single vowel letters (a, e, i, o, u, w, y). As a result, some "long vowels" are often indicated by combinations of letters (like the oa in boat, the ow in how, and the ay in stay), or the historically based silent e (as in note and cake).[243]
+
+The consequence of this complex orthographic history is that learning to read can be challenging in English. It can take longer for school pupils to become independently fluent readers of English than of many other languages, including Italian, Spanish, and German.[247] Nonetheless, there is an advantage for learners of English reading in learning the specific sound-symbol regularities that occur in the standard English spellings of commonly used words.[242] Such instruction greatly reduces the risk of children experiencing reading difficulties in English.[248][249] Making primary school teachers more aware of the primacy of morpheme representation in English may help learners learn more efficiently to read and write English.[250]
+
+English writing also includes a system of punctuation marks that is similar to those used in most alphabetic languages around the world. The purpose of punctuation is to mark meaningful grammatical relationships in sentences to aid readers in understanding a text and to indicate features important for reading a text aloud.[251]
+
+Dialectologists identify many English dialects, which usually refer to regional varieties that differ from each other in terms of patterns of grammar, vocabulary, and pronunciation. The pronunciation of particular areas distinguishes dialects as separate regional accents. The major native dialects of English are often divided by linguists into the two extremely general categories of British English (BrE) and North American English (NAE).[252] There also exists a third common major grouping of English varieties: Southern Hemisphere English, the most prominent being Australian and New Zealand English.
+
+As the place where English first evolved, the British Isles, and particularly England, are home to the most diverse dialects. Within the United Kingdom, the Received Pronunciation (RP), an educated dialect of South East England, is traditionally used as the broadcast standard and is considered the most prestigious of the British dialects. The spread of RP (also known as BBC English) through the media has caused many traditional dialects of rural England to recede, as youths adopt the traits of the prestige variety instead of traits from local dialects. At the time of the Survey of English Dialects, grammar and vocabulary differed across the country, but a process of lexical attrition has led most of this variation to disappear.[253]
+
+Nonetheless, this attrition has mostly affected dialectal variation in grammar and vocabulary, and in fact, only 3 percent of the English population actually speak RP, the remainder speaking in regional accents and dialects with varying degrees of RP influence.[254] There is also variability within RP, particularly along class lines between Upper and Middle-class RP speakers and between native RP speakers and speakers who adopt RP later in life.[255] Within Britain, there is also considerable variation along lines of social class, and some traits though exceedingly common are considered "non-standard" and are associated with lower class speakers and identities. An example of this is H-dropping, which was historically a feature of lower-class London English, particularly Cockney, and can now be heard in the local accents of most parts of England—yet it remains largely absent in broadcasting and among the upper crust of British society.[256]
+
+English in England can be divided into four major dialect regions, Southwest English, South East English, Midlands English, and Northern English. Within each of these regions several local subdialects exist: Within the Northern region, there is a division between the Yorkshire dialects and the Geordie dialect spoken in Northumbria around Newcastle, and the Lancashire dialects with local urban dialects in Liverpool (Scouse) and Manchester (Mancunian). Having been the centre of Danish occupation during the Viking Invasions, Northern English dialects, particularly the Yorkshire dialect, retain Norse features not found in other English varieties.[257]
+
+Since the 15th century, southeastern England varieties have centred on London, which has been the centre from which dialectal innovations have spread to other dialects. In London, the Cockney dialect was traditionally used by the lower classes, and it was long a socially stigmatised variety. The spread of Cockney features across the south-east led the media to talk of Estuary English as a new dialect, but the notion was criticised by many linguists on the grounds that London had been influencing neighbouring regions throughout history.[258][259][260] Traits that have spread from London in recent decades include the use of intrusive R (drawing is pronounced drawring /ˈdrɔːrɪŋ/), t-glottalisation (Potter is pronounced with a glottal stop as Po'er /poʔʌ/), and the pronunciation of th- as /f/ (thanks pronounced fanks) or /v/ (bother pronounced bover).[261]
+
+Scots is today considered a separate language from English, but it has its origins in early Northern Middle English[262] and developed and changed during its history with influence from other sources, particularly Scots Gaelic and Old Norse. Scots itself has a number of regional dialects. And in addition to Scots, Scottish English comprises the varieties of Standard English spoken in Scotland; most varieties are Northern English accents, with some influence from Scots.[263]
+
+In Ireland, various forms of English have been spoken since the Norman invasions of the 11th century. In County Wexford, in the area surrounding Dublin, two extinct dialects known as Forth and Bargy and Fingallian developed as offshoots from Early Middle English, and were spoken until the 19th century. Modern Irish English, however, has its roots in English colonisation in the 17th century. Today Irish English is divided into Ulster English, the Northern Ireland dialect with strong influence from Scots, and various dialects of the Republic of Ireland. Like Scottish and most North American accents, almost all Irish accents preserve the rhoticity which has been lost in the dialects influenced by RP.[21][264]
+
+North American English is fairly homogeneous compared to British English. Today, American accent variation is often increasing at the regional level and decreasing at the very local level,[265] though most Americans still speak within a phonological continuum of similar accents,[266] known collectively as General American (GA), with differences hardly noticed even among Americans themselves (such as Midland and Western American English).[267][268][269] In most American and Canadian English dialects, rhoticity (or r-fulness) is dominant, with non-rhoticity (r-dropping) becoming associated with lower prestige and social class especially after World War II; this contrasts with the situation in England, where non-rhoticity has become the standard.[270]
+
+Separate from GA are American dialects with clearly distinct sound systems, historically including Southern American English, English of the coastal Northeast (famously including Eastern New England English and New York City English), and African American Vernacular English, all of which are historically non-rhotic. Canadian English, except for the Atlantic provinces and perhaps Quebec, may be classified under GA as well, but it often shows the raising of the vowels /aɪ/ and /aʊ/ before voiceless consonants, as well as distinct norms for written and pronunciation standards.[271]
+
+In Southern American English, the most populous American "accent group" outside of GA,[272] rhoticity now strongly prevails, replacing the region's historical non-rhotic prestige.[273][274][275] Southern accents are colloquially described as a "drawl" or "twang,"[276] being recognised most readily by the Southern Vowel Shift initiated by glide-deleting in the /aɪ/ vowel (e.g. pronouncing spy almost like spa), the "Southern breaking" of several front pure vowels into a gliding vowel or even two syllables (e.g. pronouncing the word "press" almost like "pray-us"),[277] the pin–pen merger, and other distinctive phonological, grammatical, and lexical features, many of which are actually recent developments of the 19th century or later.[278]
+
+Today spoken primarily by working- and middle-class African Americans, African-American Vernacular English (AAVE) is also largely non-rhotic and likely originated among enslaved Africans and African Americans influenced primarily by the non-rhotic, non-standard older Southern dialects. A minority of linguists,[279] contrarily, propose that AAVE mostly traces back to African languages spoken by the slaves who had to develop a pidgin or Creole English to communicate with slaves of other ethnic and linguistic origins.[280] AAVE's important commonalities with Southern accents suggests it developed into a highly coherent and homogeneous variety in the 19th or early 20th century. AAVE is commonly stigmatised in North America as a form of "broken" or "uneducated" English, as are white Southern accents, but linguists today recognise both as fully developed varieties of English with their own norms shared by a large speech community.[281][282]
+
+Since 1788, English has been spoken in Oceania, and Australian English has developed as a first language of the vast majority of the inhabitants of the Australian continent, its standard accent being General Australian. The English of neighbouring New Zealand has to a lesser degree become an influential standard variety of the language.[283] Australian and New Zealand English are each other's closest relatives with few differentiating characteristics, followed by South African English and the English of southeastern England, all of which have similarly non-rhotic accents, aside from some accents in the South Island of New Zealand. Australian and New Zealand English stand out for their innovative vowels: many short vowels are fronted or raised, whereas many long vowels have diphthongised. Australian English also has a contrast between long and short vowels, not found in most other varieties. Australian English grammar aligns closely to British and American English; like American English, collective plural subjects take on a singular verb (as in the government is rather than are).[284][285] New Zealand English uses front vowels that are often even higher than in Australian English.[286][287][288]
+
+The first significant exposure of the Philippines to the English language occurred in 1762 when the British occupied Manila during the Seven Years' War, but this was a brief episode that had no lasting influence. English later became more important and widespread during American rule between 1898 and 1946, and remains an official language of the Philippines. Today, the use of English is ubiquitous in the Philippines, from street signs and marquees, government documents and forms, courtrooms, the media and entertainment industries, the business sector, and other aspects of daily life. One such usage that is also prominent in the country is in speech, where most Filipinos from Manila would use or have been exposed to Taglish, a form of code-switching between Tagalog and English. A similar code-switching method is used by urban native speakers of Visayan languages called Bislish.
+
+English is spoken widely in southern Africa and is an official or co-official language in several countries. In South Africa, English has been spoken since 1820, co-existing with Afrikaans and various African languages such as the Khoe and Bantu languages. Today, about 9 percent of the South African population speaks South African English (SAE) as a first language. SAE is a non-rhotic variety, which tends to follow RP as a norm. It is alone among non-rhotic varieties in lacking intrusive r. There are different L2 varieties that differ based on the native language of the speakers.[289] Most phonological differences from RP are in the vowels.[290] Consonant differences include the tendency to pronounce /p, t, t͡ʃ, k/ without aspiration (e.g. pin pronounced [pɪn] rather than as [pʰɪn] as in most other varieties), while r is often pronounced as a flap [ɾ] instead of as the more common fricative.[291]
+
+Nigerian English is a dialect of English spoken in Nigeria.[292] It is based on British English, but in recent years, because of influence from the United States, some words of American English origin have made it into Nigerian English. Additionally, some new words and collocations have emerged from the language, which come from the need to express concepts specific to the culture of the nation (e.g. senior wife). Over 150 million Nigerians speak English.[293]
+
+Several varieties of English are also spoken in the Caribbean islands that were colonial possessions of Britain, including Jamaica, and the Leeward and Windward Islands and Trinidad and Tobago, Barbados, the Cayman Islands, and Belize. Each of these areas is home both to a local variety of English and a local English based creole, combining English and African languages. The most prominent varieties are Jamaican English and Jamaican Creole. In Central America, English based creoles are spoken in on the Caribbean coasts of Nicaragua and Panama.[294] Locals are often fluent both in the local English variety and the local creole languages and code-switching between them is frequent, indeed another way to conceptualise the relationship between Creole and Standard varieties is to see a spectrum of social registers with the Creole forms serving as "basilect" and the more RP-like forms serving as the "acrolect", the most formal register.[295]
+
+Most Caribbean varieties are based on British English and consequently, most are non-rhotic, except for formal styles of Jamaican English which are often rhotic. Jamaican English differs from RP in its vowel inventory, which has a distinction between long and short vowels rather than tense and lax vowels as in Standard English. The diphthongs /ei/ and /ou/ are monophthongs [eː] and [oː] or even the reverse diphthongs [ie] and [uo] (e.g. bay and boat pronounced [bʲeː] and [bʷoːt]). Often word-final consonant clusters are simplified so that "child" is pronounced [t͡ʃail] and "wind" [win].[296][297][298]
+
+As a historical legacy, Indian English tends to take RP as its ideal, and how well this ideal is realised in an individual's speech reflects class distinctions among Indian English speakers. Indian English accents are marked by the pronunciation of phonemes such as /t/ and /d/ (often pronounced with retroflex articulation as [ʈ] and [ɖ]) and the replacement of /θ/ and /ð/ with dentals [t̪] and [d̪]. Sometimes Indian English speakers may also use spelling based pronunciations where the silent ⟨h⟩ found in words such as ghost is pronounced as an Indian voiced aspirated stop [ɡʱ].[299]
+
+Click on a coloured area to see an article about English in that country or region

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+Enid Mary Blyton (11 August 1897 – 28 November 1968) was an English children's writer whose books have been among the world's best-sellers since the 1930s, selling more than 600 million copies. Enid's books are still enormously popular, and have been translated into 90 languages. She wrote on a wide range of topics including education, natural history, fantasy, mystery, and biblical narratives and is best remembered today for her Noddy, Famous Five, Malory Towers and Secret Seven series.
+
+Her first book, Child Whispers, a 24-page collection of poems, was published in 1922. Following the commercial success of her early novels such as Adventures of the Wishing-Chair (1937) and The Enchanted Wood (1939), Blyton went on to build a literary empire, sometimes producing fifty books a year in addition to her prolific magazine and newspaper contributions. Her writing was unplanned and sprang largely from her unconscious mind: she typed her stories as events unfolded before her. The sheer volume of her work and the speed with which it was produced led to rumours that Blyton employed an army of ghost writers, a charge she vigorously denied.
+
+Blyton's work became increasingly controversial among literary critics, teachers and parents from the 1950s onwards, because of the alleged unchallenging nature of her writing and the themes of her books, particularly the Noddy series. Some libraries and schools banned her works, which the BBC had refused to broadcast from the 1930s until the 1950s because they were perceived to lack literary merit. Her books have been criticised as being elitist, sexist, racist, xenophobic and at odds with the more progressive environment emerging in post-Second World War Britain, but they have continued to be best-sellers since her death in 1968.
+
+She felt she had a responsibility to provide her readers with a strong moral framework, so she encouraged them to support worthy causes. In particular, through the clubs she set up or supported, she encouraged and organised them to raise funds for animal and paediatric charities. The story of Blyton's life was dramatised in a BBC film entitled Enid, featuring Helena Bonham Carter in the title role and first broadcast in the United Kingdom on BBC Four in 2009. There have also been several adaptations of her books for stage, screen and television.
+
+Enid Blyton was born on 11 August 1897 in East Dulwich, South London,United Kingdom the oldest of the three children, to Thomas Carey Blyton (1870–1920), a cutlery salesman, (the 1911 census records his occupation as Mantle Manufacturer dealer, womens suits, skirts, etc) and his wife Theresa Mary (née Harrison; 1874–1950). Enid's younger brothers, Hanly (1899–1983) and Carey (1902–1976), were born after the family had moved to a semi-detached villa in Beckenham, then a village in Kent.[1] A few months after her birth Enid almost died from whooping cough, but was nursed back to health by her father, whom she adored.[2] Thomas Blyton ignited Enid's interest in nature; in her autobiography she wrote that he "loved flowers and birds and wild animals, and knew more about them than anyone I had ever met".[3] He also passed on his interest in gardening, art, music, literature and the theatre, and the pair often went on nature walks, much to the disapproval of Enid's mother, who showed little interest in her daughter's pursuits.[4] Enid was devastated when she left the family shortly after her thirteenth birthday to live with another woman. Enid and her mother did not have a good relationship, and she did not attend either of her parents' funerals.[5]
+
+From 1907 to 1915 Blyton attended St Christopher's School in Beckenham, where she enjoyed physical activities and became school tennis champion and captain of lacrosse.[6] She was not so keen on all the academic subjects but excelled in writing, and in 1911 she entered Arthur Mee's children's poetry competition. Mee offered to print her verses, encouraging her to produce more.[1] Blyton's mother considered her efforts at writing to be a "waste of time and money", but she was encouraged to persevere by Mabel Attenborough, the aunt of school friend Mary Potter.[4]
+
+Blyton's father taught her to play the piano, which she mastered well enough for him to believe that she might follow in his sister's footsteps and become a professional musician.[6] Blyton considered enrolling at the Guildhall School of Music, but decided she was better suited to becoming a writer.[7] After finishing school in 1915 as head girl, she moved out of the family home to live with her friend Mary Attenborough, before going to stay with George and Emily Hunt at Seckford Hall near Woodbridge in Suffolk. Seckford Hall, with its allegedly haunted room and secret passageway provided inspiration for her later writing.[1] At Woodbridge Congregational Church Blyton met Ida Hunt, who taught at Ipswich High School, and suggested that she train there as a teacher.[8][9] Blyton was introduced to the children at the nursery school, and recognising her natural affinity with them she enrolled in a National Froebel Union teacher training course at the school in September 1916.[7][10] By this time she had almost ceased contact with her family.[1]
+
+Blyton's manuscripts had been rejected by publishers on many occasions, which only made her more determined to succeed: "it is partly the struggle that helps you so much, that gives you determination, character, self-reliance – all things that help in any profession or trade, and most certainly in writing". In March 1916 her first poems were published in Nash's Magazine.[11] She completed her teacher training course in December 1918, and the following month obtained a teaching appointment at Bickley Park School, a small independent establishment for boys in Bickley, Kent. Two months later Blyton received a teaching certificate with distinctions in zoology and principles of education, 1st class in botany, geography, practice and history of education, child hygiene and class teaching and 2nd class in literature and elementary mathematics.[1] In 1920 she moved to Southernhay in Hook Road Surbiton as nursery governess to the four sons of architect Horace Thompson and his wife Gertrude,[7] with whom Blyton spent four happy years. Owing to a shortage of schools in the area her charges were soon joined by the children of neighbours, and a small school developed at the house.[12]
+
+In 1920 Blyton relocated to Chessington, and began writing in her spare time. The following year she won the Saturday Westminster Review writing competition with her essay "On the Popular Fallacy that to the Pure All Things are Pure".[13] Publications such as The Londoner, Home Weekly and The Bystander began to show an interest in her short stories and poems.[1]
+
+Blyton's first book, Child Whispers, a 24-page collection of poems, was published in 1922.[13] It was illustrated by a schoolfriend, Phyllis Chase, who collaborated on several of her early works.[14] Also in that year Blyton began writing in annuals for Cassell and George Newnes, and her first piece of writing, "Peronel and his Pot of Glue", was accepted for publication in Teachers' World. Her success was boosted in 1923 when her poems were published alongside those of Rudyard Kipling, Walter de la Mare and G. K. Chesterton in a special issue of Teachers' World. Blyton's educational texts were quite influential in the 1920s and '30s, her most sizeable being the three-volume The Teacher's Treasury (1926), the six-volume Modern Teaching (1928), the ten-volume Pictorial Knowledge (1930), and the four-volume Modern Teaching in the Infant School (1932).[15]
+
+In July 1923 Blyton published Real Fairies, a collection of thirty-three poems written especially for the book with the exception of "Pretending", which had appeared earlier in Punch magazine.[16] The following year she published The Enid Blyton Book of Fairies, illustrated by Horace J. Knowles,[17] and in 1926 the Book of Brownies.[18] Several books of plays appeared in 1927, including A Book of Little Plays and The Play's the Thing with the illustrator Alfred Bestall.[19]
+
+In the 1930s Blyton developed an interest in writing stories related to various myths, including those of ancient Greece and Rome; The Knights of the Round Table, Tales of Ancient Greece and Tales of Robin Hood were published in 1930. In Tales of Ancient Greece Blyton retold sixteen well-known ancient Greek myths, but used the Latin rather than the Greek names of deities and invented conversations between the characters.[20] The Adventures of Odysseus, Tales of the Ancient Greeks and Persians and Tales of the Romans followed in 1934.[21]
+
+The first of twenty-eight books in Blyton's Old Thatch series, The Talking Teapot and Other Tales, was published in 1934, the same year as Brer Rabbit Retold;[22] (note that Brer Rabbit originally featured in Uncle Remus stories by Joel Chandler Harris), her first serial story and first full-length book, Adventures of the Wishing-Chair, followed in 1937. The Enchanted Wood, the first book in the Faraway Tree series, published in 1939, is about a magic tree inspired by the Norse mythology that had fascinated Blyton as a child.[7] According to Blyton's daughter Gillian the inspiration for the magic tree came from "thinking up a story one day and suddenly she was walking in the enchanted wood and found the tree. In her imagination she climbed up through the branches and met Moon-Face, Silky, the Saucepan Man and the rest of the characters. She had all she needed."[23] As in the Wishing-Chair series, these fantasy books typically involve children being transported into a magical world in which they meet fairies, goblins, elves, pixies and other mythological creatures.
+
+Blyton's first full-length adventure novel, The Secret Island, was published in 1938, featuring the characters of Jack, Mike, Peggy and Nora.[24] Described by The Glasgow Herald as a "Robinson Crusoe-style adventure on an island in an English lake", The Secret Island was a lifelong favourite of Gillian's and spawned the Secret series.[23] The following year Blyton released her first book in the Circus series[25] and her initial book in the Amelia Jane series, Naughty Amelia Jane![26] According to Gillian the main character was based on a large handmade doll given to her by her mother on her third birthday.[23]
+
+During the 1940s Blyton became a prolific author, her success enhanced by her "marketing, publicity and branding that was far ahead of its time".[27] In 1940 Blyton published two books – Three Boys and a Circus and Children of Kidillin – under the pseudonym of Mary Pollock (middle name plus first married name),[28] in addition to the eleven published under her own name that year. So popular were Pollock's books that one reviewer was prompted to observe that "Enid Blyton had better look to her laurels".[29] But Blyton's readers were not so easily deceived and many complained about the subterfuge to her and her publisher,[29] with the result that all six books published under the name of Mary Pollock – two in 1940 and four in 1943 – were reissued under Blyton's name.[30] Later in 1940 Blyton published the first of her boarding school story books and the first novel in the Naughtiest Girl series, The Naughtiest Girl in the School, which followed the exploits of the mischievous schoolgirl Elizabeth Allen at the fictional Whyteleafe School. The first of her six novels in the St. Clare's series, The Twins at St. Clare's, appeared the following year, featuring the twin sisters Patricia and Isabel O'Sullivan.[15]
+
+In 1942 Blyton released the first book in the Mary Mouse series, Mary Mouse and the Dolls' House, about a mouse exiled from her mousehole who becomes a maid at a dolls' house. Twenty-three books in the series were produced between 1942 and 1964; 10,000 copies were sold in 1942 alone.[31] The same year, Blyton published the first novel in the Famous Five series, Five on a Treasure Island, with illustrations by Eileen Soper. Its popularity resulted in twenty-one books between then and 1963, and the characters of Julian, Dick, Anne, George (Georgina) and Timmy the dog became household names in Britain.[32] Matthew Grenby, author of Children's Literature, states that the five were involved with "unmasking hardened villains and solving serious crimes", although the novels were "hardly 'hard-boiled' thrillers".[33] Blyton based the character of Georgina, a tomboy she described as "short-haired, freckled, sturdy, and snub-nosed" and "bold and daring, hot-tempered and loyal", on herself.[11]
+
+Blyton had an interest in biblical narratives, and retold Old and New Testament stories. The Land of Far-Beyond (1942) is a Christian parable along the lines of John Bunyan's The Pilgrim's Progress (1698), with contemporary children as the main characters.[34] In 1943 she published The Children's Life of Christ, a collection of fifty-nine short stories related to the life of Jesus, with her own slant on popular biblical stories, from the Nativity and the Three Wise Men through to the trial, the crucifixion and the resurrection.[35] Tales from the Bible was published the following year,[36] followed by The Boy with the Loaves and Fishes in 1948.[37]
+
+The first book in Blyton's Five Find-Outers series, The Mystery of the Burnt Cottage, was published in 1943, as was the second book in the Faraway series, The Magic Faraway Tree, which in 2003 was voted 66th in the BBC's Big Read poll to find the UK's favourite book.[38] Several of Blyton's works during this period have seaside themes; John Jolly by the Sea (1943), a picture book intended for younger readers, was published in a booklet format by Evans Brothers.[39] Other books with a maritime theme include The Secret of Cliff Castle and Smuggler Ben, both attributed to Mary Pollock in 1943;[40] The Island of Adventure, the first in the Adventure series of eight novels from 1944 onwards;[41] and various novels of the Famous Five series such as Five on a Treasure Island (1942),[42] Five on Kirrin Island Again (1947)[43] and Five Go Down to the Sea (1953).[44]
+
+Capitalising on her success, with a loyal and ever-growing readership,[15] Blyton produced a new edition of many of her series such as the Famous Five, the Five Find-Outers and St. Clare's every year in addition to many other novels, short stories and books. In 1946 Blyton launched the first in the Malory Towers series of six books based around the schoolgirl Darrell Rivers, First Term at Malory Towers, which became extremely popular, particularly with girls.[45]
+
+The first book in Blyton's Barney Mysteries series, The Rockingdown Mystery, was published in 1949,[46] as was the first of her fifteen Secret Seven novels.[47] The Secret Seven Society consists of Peter, his sister Janet, and their friends Colin, George, Jack, Pam and Barbara, who meet regularly in a shed in the garden to discuss peculiar events in their local community. Blyton rewrote the stories so they could be adapted into cartoons, which appeared in Mickey Mouse Weekly in 1951 with illustrations by George Brook. The French author Evelyne Lallemand continued the series in the 1970s, producing an additional twelve books, nine of which were translated into English by Anthea Bell between 1983 and 1987.[48]
+
+Blyton's Noddy, about a little wooden boy from Toyland, first appeared in the Sunday Graphic on 5 June 1949, and in November that year Noddy Goes to Toyland, the first of at least two dozen books in the series, was published. The idea was conceived by one of Blyton's publishers, Sampson, Low, Marston and Company, who in 1949 arranged a meeting between Blyton and the Dutch illustrator Harmsen van der Beek. Despite having to communicate via an interpreter, he provided some initial sketches of how Toyland and its characters would be represented. Four days after the meeting Blyton sent the text of the first two Noddy books to her publisher, to be forwarded to van der Beek.[49] The Noddy books became one of her most successful and best-known series, and were hugely popular in the 1950s.[50] An extensive range of sub-series, spin-offs and strip books were produced throughout the decade, including Noddy's Library, Noddy's Garage of Books, Noddy's Castle of Books, Noddy's Toy Station of Books and Noddy's Shop of Books.[51]
+
+In 1950 Blyton established the company Darrell Waters Ltd to manage her affairs. By the early 1950s she had reached the peak of her output, often publishing more than fifty books a year, and she remained extremely prolific throughout much of the decade.[52] By 1955 Blyton had written her fourteenth Famous Five novel, Five Have Plenty of Fun, her fifteenth Mary Mouse book, Mary Mouse in Nursery Rhyme Land, her eighth book in the Adventure series, The River of Adventure, and her seventh Secret Seven novel, Secret Seven Win Through. She completed the sixth and final book of the Malory Towers series, Last Term at Malory Towers, in 1951.[45]
+
+Blyton published several further books featuring the character of Scamp the terrier, following on from The Adventures of Scamp, a novel she had released in 1943 under the pseudonym of Mary Pollock.[53] Scamp Goes on Holiday (1952) and Scamp and Bimbo, Scamp at School, Scamp and Caroline and Scamp Goes to the Zoo (1954) were illustrated by Pierre Probst. She introduced the character of Bom, a stylish toy drummer dressed in a bright red coat and helmet, alongside Noddy in TV Comic in July 1956.[54] A book series began the same year with  Bom the Little Toy Drummer, featuring illustrations by R. Paul-Hoye,[55] and followed with Bom and His Magic Drumstick (1957), Bom Goes Adventuring and Bom Goes to Ho Ho Village (1958), Bom and the Clown and Bom and the Rainbow (1959) and Bom Goes to Magic Town (1960). In 1958 she produced two annuals featuring the character, the first of which included twenty short stories, poems and picture strips.[56]
+
+Many of Blyton's series, including Noddy and The Famous Five, continued to be successful in the 1960s; by 1962, 26 million copies of Noddy had been sold.[1][a] Blyton concluded several of her long-running series in 1963, publishing the last books of The Famous Five (Five Are Together Again) and The Secret Seven (Fun for the Secret Seven); she also produced three more Brer Rabbit books with the illustrator Grace Lodge: Brer Rabbit Again, Brer Rabbit Book, and Brer Rabbit's a Rascal. In 1962 many of her books were among the first to be published by Armada Books in paperback, making them more affordable to children.[1]
+
+After 1963 Blyton's output was generally confined to short stories and books intended for very young readers, such as Learn to Count with Noddy and Learn to Tell Time with Noddy in 1965, and Stories for Bedtime and the Sunshine Picture Story Book collection in 1966. Her declining health and a falling off in readership among older children have been put forward as the principal reasons for this change in trend.[57] Blyton published her last book in the Noddy series, Noddy and the Aeroplane, in February 1964. In May the following year she published Mixed Bag, a song book with music written by her nephew Carey, and in August she released her last full-length books, The Man Who Stopped to Help and The Boy Who Came Back.[1]
+
+Blyton cemented her reputation as a children's writer when in 1926 she took over the editing of Sunny Stories, a magazine that typically included the re-telling of legends, myths, stories and other articles for children.[7] That same year she was given her own column in Teachers' World, entitled "From my Window". Three years later she began contributing a weekly page in the magazine, in which she published letters from her fox terrier dog Bobs.[1] They proved to be so popular that in 1933 they were published in book form as Letters from Bobs,[58] and sold ten thousand copies in the first week.[1] Her most popular feature was "Round the Year with Enid Blyton", which consisted of forty-eight articles covering aspects of natural history such as weather, pond life, how to plant a school garden and how to make a bird table.[59] Among Blyton's other nature projects was her monthly "Country Letter" feature that appeared in The Nature Lover magazine in 1935.[60]
+
+Sunny Stories was renamed Enid Blyton's Sunny Stories in January 1937, and served as a vehicle for the serialisation of Blyton's books. Her first Naughty Amelia Jane story, about an anti-heroine based on a doll owned by her daughter Gillian,[61] was published in the magazine.[1] Blyton stopped contributing in 1952, and it closed down the following year, shortly before the appearance of the new fortnightly Enid Blyton Magazine written entirely by Blyton.[62] The first edition appeared on 18 March 1953,[63] and the magazine ran until September 1959.[7]
+
+Noddy made his first appearance in the Sunday Graphic in 1949, the same year as Blyton's first daily Noddy strip for the London Evening Standard.[1] It was illustrated by van der Beek until his death in 1953.[1][64]
+
+Blyton worked in a wide range of fictional genres, from fairy tales to animal, nature, detective, mystery, and circus stories, but she often "blurred the boundaries" in her books, and encompassed a range of genres even in her short stories.[65] In a 1958 article published in The Author, she wrote that there were a "dozen or more different types of stories for children", and she had tried them all, but her favourites were those with a family at their centre.[66]
+
+In a letter to the psychologist Peter McKellar,[b] Blyton describes her writing technique:
+
+I shut my eyes for a few minutes, with my portable typewriter on my knee – I make my mind a blank and wait – and then, as clearly as I would see real children, my characters stand before me in my mind's eye ... The first sentence comes straight into my mind, I don't have to think of it – I don't have to think of anything.[68]
+
+In another letter to McKellar she describes how in just five days she wrote the 60,000-word book The River of Adventure, the eighth in her Adventure Series,[69] by listening to what she referred to as her "under-mind",[70] which she contrasted with her "upper conscious mind".[71] Blyton was unwilling to conduct any research or planning before beginning work on a new book, which coupled with the lack of variety in her life[c] according to Druce almost inevitably presented the danger that she might unconsciously, and clearly did, plagiarise the books she had read, including her own.[72] Gillian has recalled that her mother "never knew where her stories came from", but that she used to talk about them "coming from her 'mind's eye'", as did William Wordsworth and Charles Dickens. Blyton had "thought it was made up of every experience she'd ever had, everything she's seen or heard or read, much of which had long disappeared from her conscious memory" but never knew the direction her stories would take.  Blyton further explained in her biography that "If I tried to think out or invent the whole book, I could not do it. For one thing, it would bore me and for another, it would lack the 'verve' and the extraordinary touches and surprising ideas that flood out from my imagination."[23]
+
+Blyton's daily routine varied little over the years. She usually began writing soon after breakfast, with her portable typewriter on her knee and her favourite red Moroccan shawl nearby; she believed that the colour red acted as a "mental stimulus" for her. Stopping only for a short lunch break she continued writing until five o'clock, by which time she would usually have produced 6,000–10,000 words.[74]
+
+A 2000 article in The Malay Mail considers Blyton's children to have "lived in a world shaped by the realities of post-war austerity", enjoying freedom without the political correctness of today, which serves modern readers of Blyton's novels with a form of escapism.[75] Brandon Robshaw of The Independent refers to the Blyton universe as  "crammed with colour and character", "self-contained and internally consistent", noting that Blyton exemplifies a strong mistrust of adults and figures of authority in her works, creating a world in which children govern.[76] Gillian noted that in her mother's adventure, detective and school stories for older children, "the hook is the strong storyline with plenty of cliffhangers, a trick she acquired from her years of writing serialised stories for children's magazines. There is always a strong moral framework in which bravery and loyalty are (eventually) rewarded".[23] Blyton herself wrote that "my love of children is the whole foundation of all my work".[77]
+
+Victor Watson, Assistant Director of Research at Homerton College, Cambridge, believes that Blyton's works reveal an "essential longing and potential associated with childhood", and notes how the opening pages of The Mountain of Adventure present a "deeply appealing ideal of childhood".[78] He argues that Blyton's work differs from that of many other authors in its approach, describing the narrative of The Famous Five series for instance as "like a powerful spotlight, it seeks to illuminate, to explain, to demystify. It takes its readers on a roller-coaster story in which the darkness is always banished; everything puzzling, arbitrary, evocative is either dismissed or explained". Watson further notes how Blyton often used minimalist visual descriptions and introduced a few careless phrases such as "gleamed enchantingly" to appeal to her young readers.[79]
+
+From the mid-1950s rumours began to circulate that Blyton had not written all the books attributed to her, a charge she found particularly distressing. She published an appeal in her magazine asking children to let her know if they heard such stories and, after one mother informed her that she had attended a parents' meeting at her daughter's school during which a young librarian had repeated the allegation,[80] Blyton decided in 1955 to begin legal proceedings.[1] The librarian was eventually forced to make a public apology in open court early the following year, but the rumours that Blyton operated "a 'company' of ghost writers" persisted, as some found it difficult to believe that one woman working alone could produce such a volume of work.[81]
+
+Enid's Conservative personal politics were often in view in her fiction. In The Mystery of the Missing Necklace (a The Five Find-Outers installment), she uses the character of young Elizabeth ("Bets") to give a statement praising Winston Churchill and describing the politician as a "statesman".[82]
+
+Blyton felt a responsibility to provide her readers with a positive moral framework, and she encouraged them to support worthy causes.[83] Her view, expressed in a 1957 article, was that children should help animals and other children rather than adults:
+
+[children] are not interested in helping adults; indeed, they think that adults themselves should tackle adult needs. But they are intensely interested in animals and other children and feel compassion for the blind boys and girls, and for the spastics who are unable to walk or talk.[84]
+
+Blyton and the members of the children's clubs she promoted via her magazines raised a great deal of money for various charities; according to Blyton, membership of her clubs meant "working for others, for no reward". The largest of the clubs she was involved with was the Busy Bees, the junior section of the People's Dispensary for Sick Animals, which Blyton had actively supported since 1933. The club had been set up by Maria Dickin in 1934,[85] and after Blyton publicised its existence in the Enid Blyton Magazine it attracted 100,000 members in three years.[86] Such was Blyton's popularity among children that after she became Queen Bee in 1952 more than 20,000 additional members were recruited in her first year in office.[85] The Enid Blyton Magazine Club was formed in 1953.[1] Its primary objective was to raise funds to help those children with cerebral palsy who attended a centre in Cheyne Walk, in Chelsea, London, by furnishing an on-site hostel among other things.[87]
+
+The Famous Five series gathered such a following that readers asked Blyton if they might form a fan club. She agreed, on condition that it serve a useful purpose, and suggested that it could raise funds for the Shaftesbury Society Babies' Home[d] in Beaconsfield, on whose committee she had served since 1948.[89] The club was established in 1952, and provided funds for equipping a Famous Five Ward at the home, a paddling pool, sun room, summer house, playground, birthday and Christmas celebrations, and visits to the pantomime.[88] By the late 1950s Blyton's clubs had a membership of 500,000, and raised £35,000 in the six years of the Enid Blyton Magazine's run.[4]
+
+By 1974 the Famous Five Club had a membership of 220,000, and was growing at the rate of 6,000 new members a year.[90][e] The Beaconsfield home it was set up to support closed in 1967, but the club continued to raise funds for other paediatric charities, including an Enid Blyton bed at Great Ormond Street Hospital and a mini-bus for disabled children at Stoke Mandeville Hospital.[92]
+
+Blyton capitalised upon her commercial success as an author by negotiating agreements with jigsaw puzzle and games manufacturers from the late 1940s onwards; by the early 1960s some 146 different companies were involved in merchandising Noddy alone.[93] In 1948 Bestime released four jigsaw puzzles featuring her characters, and the first Enid Blyton board game appeared, Journey Through Fairyland, created by BGL. The first card game, Faraway Tree, appeared from Pepys in 1950. In 1954 Bestime released the first four jigsaw puzzles of the Secret Seven, and the following year a Secret Seven card game appeared.[48]
+
+Bestime released the Little Noddy Car Game in 1953 and the Little Noddy Leap Frog Game in 1955, and in 1956 American manufacturer Parker Brothers released Little Noddy's Taxi Game, a board game which features Noddy driving about town, picking up various characters.[94] Bestime released its Plywood Noddy Jigsaws series in 1957 and a Noddy jigsaw series featuring cards appeared from 1963, with illustrations by Robert Lee. Arrow Games became the chief producer of Noddy jigsaws in the late 1970s and early 1980s.[93] Whitman manufactured four new Secret Seven jigsaw puzzles in 1975, and produced four new Malory Towers ones two years later.[48] In 1979 the company released a Famous Five adventure board game, Famous Five Kirrin Island Treasure.[95] Stephen Thraves wrote eight Famous Five adventure game books, published by Hodder & Stoughton in the 1980s. The first adventure game book of the series, The Wreckers' Tower Game, was published in October 1984.[96]
+
+On 28 August 1924 Blyton married Major Hugh Alexander Pollock, DSO (1888–1971) at Bromley Register Office, without inviting her family.[1] They married shortly after he divorced from his first wife, with whom he had two sons, one of whom was already deceased. Pollock was editor of the book department in the publishing firm of George Newnes, which became her regular publisher. It was he who requested that Blyton write a book about animals, The Zoo Book, which was completed in the month before they married.[1] They initially lived in a flat in Chelsea before moving to Elfin Cottage in Beckenham in 1926, and then to Old Thatch in Bourne End (called Peterswood in her books) in 1929.[7][97] Blyton's first daughter Gillian, was born on 15 July 1931, and after a miscarriage in 1934,[4] she gave birth to a second daughter, Imogen, on 27 October 1935.[1]
+
+In 1938 Blyton and her family moved to a house in Beaconsfield, which was named Green Hedges by Blyton's readers following a competition in her magazine. By the mid-1930s, Pollock – possibly due to the trauma he had suffered during the First World War being revived through his meetings as a publisher with Winston Churchill – withdrew increasingly from public life and became a secret alcoholic.[98] With the outbreak of the Second World War, he became involved in the Home Guard.[98] Pollock met again Ida Crowe, an aspiring writer nineteen years his junior, whom he had met years before. He made an offer to her to join him as secretary in his posting to a Home Guard training centre at Denbies, a Gothic mansion in Surrey belonging to Lord Ashcombe, and they entered into a romantic relationship.[99] Blyton's marriage to Pollock became troubled for years, and according to Crowe's memoir, Blyton began a series of affairs,[99] including a lesbian relationship with one of the children's nannies.[99][100]  In 1941 Blyton met Kenneth Fraser Darrell Waters, a London surgeon with whom she began a serious affair.[101] Pollock discovered the liaison, and threatened to initiate divorce proceedings against Blyton.[102] Fearing that exposure of her adultery would ruin her public image,[99] it was ultimately agreed that Blyton would instead file for divorce against Pollock.[102] According to Crowe's memoir, Blyton promised that if he admitted to infidelity she would allow him parental access to their daughters; but after the divorce he was forbidden to contact them, and Blyton ensured he was subsequently unable to find work in publishing. Pollock, having married Crowe on 26 October 1943, eventually resumed his heavy drinking and was forced to petition for bankruptcy in 1950.[99]
+
+Blyton and Darrell Waters married at the City of Westminster Register Office on 20 October 1943. She changed the surname of her daughters to Darrell Waters[103] and publicly embraced her new role as a happily married and devoted doctor's wife.[7] After discovering she was pregnant in the spring of 1945, Blyton miscarried five months later, following a fall from a ladder. The baby would have been Darrell Waters's first child and it would also have been the son for which both of them longed.[4]
+
+Her love of tennis included playing naked, with nude tennis "a common practice in those days among the more louche members of the middle classes".[104][105]
+
+Blyton's health began to deteriorate in 1957, when during a round of golf she started to complain of feeling faint and breathless,[106] and by 1960 she was displaying signs of dementia.[107] Her agent George Greenfield recalled that it was "unthinkable" for the "most famous and successful of children's authors with her enormous energy and computer-like memory" to be losing her mind and suffering from what is now known as Alzheimer's disease in her mid-sixties.[107] Blyton's situation was worsened by her husband's declining health throughout the 1960s; he suffered from severe arthritis in his neck and hips, deafness, and became increasingly ill-tempered and erratic until his death on 15 September 1967.[101][108]
+
+The story of Blyton's life was dramatised in a BBC film entitled Enid, which aired in the United Kingdom on BBC Four on 16 November 2009.[109] Helena Bonham Carter, who played the title role, described Blyton as "a complete workaholic, an achievement junkie and an extremely canny businesswoman" who "knew how to brand herself, right down to the famous signature".[27]
+
+During the months following her husband's death, Blyton became increasingly ill and moved into a nursing home three months before her death. She died at the Greenways Nursing Home, Hampstead, North London, on 28 November 1968, aged 71. A memorial service was held at St James's Church, Piccadilly[1] and she was cremated at Golders Green Crematorium, where her ashes remain. Blyton's home, Green Hedges, was auctioned on 26 May 1971 and demolished in 1973;[110] the site is now occupied by houses and a street named Blyton Close. An English Heritage blue plaque commemorates Blyton at Hook Road in Chessington, where she lived from 1920 to 1924.[111] In 2014, a plaque recording her time as a Beaconsfield resident from 1938 until her death in 1968 was unveiled in the town hall gardens, next to small iron figures of Noddy and Big Ears.[112]
+
+Since her death and the publication of her daughter Imogen's 1989 autobiography, A Childhood at Green Hedges, Blyton has emerged as an emotionally immature, unstable and often malicious figure.[27] Imogen considered her mother to be "arrogant, insecure, pretentious, very skilled at putting difficult or unpleasant things out of her mind, and without a trace of maternal instinct. As a child, I viewed her as a rather strict authority. As an adult I pitied her."[113] Blyton's eldest daughter Gillian remembered her rather differently however, as "a fair and loving mother, and a fascinating companion".[113]
+
+The Enid Blyton Trust for Children was established in 1982, with Imogen as its first chairman,[114] and in 1985 it established the National Library for the Handicapped Child.[7] Enid Blyton's Adventure Magazine began publication in September 1985 and, on 14 October 1992, the BBC began publishing Noddy Magazine and released the Noddy CD-Rom in October 1996.[1]
+
+The first Enid Blyton Day was held at Rickmansworth on 6 March 1993 and, in October 1996, the Enid Blyton award, The Enid, was given to those who have made outstanding contributions towards children.[1] The Enid Blyton Society was formed in early 1995, to provide "a focal point for collectors and enthusiasts of Enid Blyton" through its thrice-annual Enid Blyton Society Journal, its annual Enid Blyton Day and its website.[115] On 16 December 1996, Channel 4 broadcast a documentary about Blyton, Secret Lives. To celebrate her centenary in 1997, exhibitions were put on at the London Toy & Model Museum (now closed), Hereford and Worcester County Museum and Bromley Library and, on 9 September, the Royal Mail issued centenary stamps.[1]
+
+The London-based entertainment and retail company Trocadero plc purchased Blyton's Darrell Waters Ltd in 1995 for £14.6 million and established a subsidiary, Enid Blyton Ltd, to handle all intellectual properties, character brands and media in Blyton's works.[1][7] The group changed its name to Chorion in 1998 but, after financial difficulties in 2012, sold its assets. Hachette UK acquired from Chorion world rights in the Blyton estate in March 2013, including The Famous Five series[116] but excluding the rights to Noddy, which had been sold to DreamWorks Classics (formerly Classic Media, now a subsidiary of DreamWorks Animation)[117] in 2012.
+
+Blyton's granddaughter, Sophie Smallwood, wrote a new Noddy book to celebrate the character's 60th birthday, 46 years after the last book was published; Noddy and the Farmyard Muddle (2009) was illustrated by Robert Tyndall.[118] In February 2011, the manuscript of a previously unknown Blyton novel, Mr Tumpy's Caravan, was discovered by the archivist at Seven Stories, National Centre for Children's Books in a collection of papers belonging to Blyton's daughter Gillian, purchased by Seven Stories in 2010 following her death.[119][120] It was initially thought to belong to a comic strip collection of the same name published in 1949, but it appears to be unrelated and is believed to be something written in the 1930s, which had been rejected by a publisher.[120][121]
+
+In a 1982 survey of 10,000 eleven-year-old children, Blyton was voted their most popular writer.[1] She is the world's fourth most-translated author, behind Agatha Christie, Jules Verne and William Shakespeare[122] with her books being translated into 90 languages.[123] From 2000 to 2010, Blyton was listed as a Top Ten author, selling almost 8 million copies (worth £31.2 million) in the UK alone.[124] In 2003, The Magic Faraway Tree was voted 66th in the BBC's Big Read.[38]  In the 2008 Costa Book Awards, Blyton was voted Britain's best-loved author.[125][126] Her books continue to be very popular among children in Commonwealth nations such as India, Pakistan, Sri Lanka, Singapore, Malta, New Zealand and Australia, and around the world.[127] They have also seen a surge of popularity in China, where they are "big with every generation".[75] In March 2004, Chorion and the Chinese publisher Foreign Language Teaching and Research Press negotiated an agreement over the Noddy franchise, which included bringing the character to an animated series on television, with a potential audience of a further 95 million children under the age of five.[128][129] Chorion spent around £10 million digitising Noddy and, as of 2002, had made television agreements with at least 11 countries worldwide.[130]
+
+Novelists influenced by Blyton include the crime writer Denise Danks, whose fictional detective Georgina Powers is based on George from the Famous Five. Peter Hunt's A Step off the Path (1985) is also influenced by the Famous Five, and the St. Clare's and Malory Towers series provided the inspiration for Jacqueline Wilson's Double Act (1996) and Adèle Geras's Egerton Hall trilogy (1990–92) respectively.[131]
+
+A.H. Thompson, who compiled an extensive overview of censorship efforts in the United Kingdom's public libraries, dedicated an entire chapter to "The Enid Blyton Affair", and wrote of her in 1975:
+
+"No single author has caused more controversy among librarians, literary critics, teachers, and other educationalists and parents during the last thirty years, than Enid Blyton. How is it that the books of this tremendously popular writer for children should have given rise to accusations of censorship against librarians in Australia, New Zealand, and the United Kingdom?"[132]
+
+Blyton's range of plots and settings has been described as limited, repetitive and continually recycled.[72] Many of her books were critically assessed by teachers and librarians, deemed unfit for children to read, and removed from syllabuses and public libraries.[7] Responding to claims that her moral views were "dependably predictable",[133] Blyton commented that "most of you could write down perfectly correctly all the things that I believe in and stand for – you have found them in my books, and a writer's books are always a faithful reflection of himself".[134]
+
+From the 1930s to the 1950s the BBC operated a de facto ban on dramatising Blyton's books for radio, considering her to be a "second-rater" whose work was without literary merit.[135][136][f] The children's literary critic Margery Fisher likened Blyton's books to "slow poison",[7] and Jean E. Sutcliffe of the BBC's schools broadcast department wrote of Blyton's ability to churn out "mediocre material", noting that "her capacity to do so amounts to genius ... anyone else would have died of boredom long ago".[137] Michael Rosen, Children's Laureate from 2007 until 2009, wrote that "I find myself flinching at occasional bursts of snobbery and the assumed level of privilege of the children and families in the books."[123] The children's author Anne Fine presented an overview of the concerns about Blyton's work and responses to them on BBC Radio 4 in November 2008, in which she noted the "drip, drip, drip of disapproval" associated with the books.[138] Blyton's response to her critics was that she was uninterested in the views of anyone over the age of 12, claiming that half the attacks on her work were motivated by jealousy and the rest came from "stupid people who don't know what they're talking about because they've never read any of my books".[139]
+
+Despite criticism by contemporaries that her work's quality began to suffer in the 1950s at the expense of its increasing volume, Blyton nevertheless capitalised on being generally regarded at the time as "a more 'savoury', English alternative" to what some considered an "invasion" of Britain by American culture, in the form of "rock music, horror comics, television, teenage culture, delinquency, and Disney".[15]
+
+According to British academic Nicholas Tucker, the works of Enid Blyton have been "banned from more public libraries over the years than is the case with any other adult or children's author", though such attempts to quell the popularity of her books over the years seem to have been largely unsuccessful, and "she still remains very widely read".[140]
+
+Some librarians felt that Blyton's restricted use of language, a conscious product of her teaching background, was prejudicial to an appreciation of more literary qualities. In a scathing article published in Encounter in 1958, the journalist Colin Welch remarked that it was "hard to see how a diet of Miss Blyton could help with the 11-plus or even with the Cambridge English Tripos",[7] but reserved his harshest criticism for Blyton's Noddy, describing him as an "unnaturally priggish ... sanctimonious ... witless, spiritless, snivelling, sneaking doll."[57]
+
+The author and educational psychologist Nicholas Tucker notes that it was common to see Blyton cited as people's favourite or least favourite author according to their age, and argues that her books create an "encapsulated world for young readers that simply dissolves with age, leaving behind only memories of excitement and strong identification".[141] Fred Inglis considers Blyton's books to be technically easy to read, but to also be "emotionally and cognitively easy". He mentions that the psychologist Michael Woods believed that Blyton was different from many other older authors writing for children in that she seemed untroubled by presenting them with a world that differed from reality. Woods surmised that Blyton "was a child, she thought as a child, and wrote as a child ... the basic feeling is essentially pre-adolescent ... Enid Blyton has no moral dilemmas ... Inevitably Enid Blyton was labelled by rumour a child-hater. If true, such a fact should come as no surprise to us, for as a child herself all other children can be nothing but rivals for her."[142] Inglis argues though that Blyton was clearly devoted to children and put an enormous amount of energy into her work, with a powerful belief in "representing the crude moral diagrams and garish fantasies of a readership".[142] Blyton's daughter Imogen has stated that she "loved a relationship with children through her books", but real children were an intrusion, and there was no room for intruders in the world that Blyton occupied through her writing.[143]
+
+Accusations of racism in Blyton's books were first made by Lena Jeger in a Guardian article published in 1966, in which she was critical of Blyton's The Little Black Doll, published a few months earlier. Sambo, the black doll of the title, is hated by his owner and the other toys owing to his "ugly black face", and runs away. A shower of rain washes his face clean, after which he is welcomed back home with his now pink face.[144] Jamaica Kincaid also considers the Noddy books to be "deeply racist" because of the blonde children and the black golliwogs.[145] In Blyton's 1944 novel The Island of Adventure, a black servant named Jo-Jo is very intelligent, but is particularly cruel to the children.[146]
+
+Accusations of xenophobia were also made. As George Greenfield observed, "Enid was very much part of that between-the-wars middle class which believed that foreigners were untrustworthy or funny or sometimes both".[147] The publisher Macmillan conducted an internal assessment of Blyton's The Mystery That Never Was, submitted to them at the height of her fame in 1960. The review was carried out by the author and books editor Phyllis Hartnoll, in whose view "There is a faint but unattractive touch of old-fashioned xenophobia in the author's attitude to the thieves; they are 'foreign' ... and this seems to be regarded as sufficient to explain their criminality." Macmillan rejected the manuscript,[148] but it was published by William Collins in 1961,[149] and then again in 1965 and 1983.[148]
+
+Blyton's depictions of boys and girls are considered by many critics to be sexist.[150][151] In a Guardian article published in 2005 Lucy Mangan proposed that The Famous Five series depicts a power struggle between Julian, Dick and George (Georgina), in which the female characters either act like boys or are talked down to, as when Dick lectures George: "it's really time you gave up thinking you're as good as a boy".[152]
+
+In December 2016 the Royal Mint discussed featuring Blyton on a commemorative 50p coin but dismissed the idea because she was "known to have been a racist, sexist, homophobe and not a very well-regarded writer".[153]
+
+To address criticisms levelled at Blyton's work some later editions have been altered to reflect more politically progressive attitudes towards issues such as race, gender, violence between young persons, the treatment of children by adults, and legal changes in Britain as to what is allowable for young children to do in the years since the stories were originally written (e.g. purchasing fireworks); modern reprints of the Noddy series substitute teddy bears or goblins for golliwogs, for instance.[154] The golliwogs who steal Noddy's car and dump him naked in the Dark Wood in Here Comes Noddy Again are replaced by goblins in the 1986 revision, who strip Noddy only of his shoes and hat and return at the end of the story to apologise.[155]
+
+The Faraway Tree's Dame Slap, who made regular use of corporal punishment, was changed to Dame Snap who no longer did so, and the names of Dick and Fanny in the same series were changed to Rick and Frannie.[156] Characters in the Malory Towers and St. Clare's series are no longer spanked or threatened with a spanking, but are instead scolded. References to George's short hair making her look like a boy were removed in revisions to Five on a Hike Together, reflecting the idea that girls need not have long hair to be considered feminine or normal.[157] Anne of The Famous Five stating that boys cannot wear pretty dresses or like girl's dolls was removed.[158] In The Adventurous Four, the names of the young twin girls were changed from Jill and Mary to Pippa and Zoe.[159]
+
+In 2010 Hodder, the publisher of the Famous Five series, announced its intention to update the language used in the books, of which it sold more than half a million copies a year. The changes, which Hodder described as "subtle", mainly affect the dialogue rather than the narrative. For instance, "school tunic" becomes "uniform", "mother and father", and "mother and daddy" (this latter one used by young female characters and deemed sexist) becomes "mum and dad",[160] "bathing" is replaced by "swimming", and "jersey" by "jumper".[156] Some commentators see the changes as necessary to encourage modern readers,[160] whereas others regard them as unnecessary and patronising.[156] In 2016 Hodder's parent company Hachette announced that they would abandon the revisions as, based on feedback, they had not been a success.[161]
+
+In 1954 Blyton adapted Noddy for the stage, producing the Noddy in Toyland pantomime in just two or three weeks. The production was staged at the 2660-seat Stoll Theatre in Kingsway, London at Christmas.[162] Its popularity resulted in the show running during the Christmas season for five or six years.[163] Blyton was delighted with its reception by children in the audience, and attended the theatre three or four times a week.[164] TV adaptations of Noddy since 1954 include one in the 1970s narrated by Richard Briers.[165] In 1955 a stage play based on the Famous Five was produced, and in January 1997 the King's Head Theatre embarked on a six-month tour of the UK with The Famous Five Musical, to commemorate Blyton's centenary. On 21 November 1998 The Secret Seven Save the World was first performed at the Sherman Theatre in Cardiff.[1]
+
+There have also been several film and television adaptations of the Famous Five: by the Children's Film Foundation in 1957 and 1964, Southern Television in 1978–79, and Zenith Productions in 1995–97.[7] The series was also adapted for the German film Fünf Freunde, directed by Mike Marzuk and released in 2011.[166]
+
+The Comic Strip, a group of British comedians, produced two extreme parodies of the Famous Five for Channel 4 television: Five Go Mad in Dorset, broadcast in 1982,[g] and Five Go Mad on Mescalin, broadcast the following year.[1] A third in the series, Five Go to Rehab, was broadcast on Sky in 2012.[167]
+
+Blyton's The Faraway Tree series of books has also been adapted to television and film. On 29 September 1997 the BBC began broadcasting an animated series called The Enchanted Lands, based on the series.[168] It was announced in October 2014 that a deal had been signed with publishers Hachette for "The Faraway Tree" series to be adapted into a live-action film by director Sam Mendes’ production company. Marlene Johnson, head of children's books at Hachette, said: "Enid Blyton was a passionate advocate of children’s storytelling, and The Magic Faraway Tree is a fantastic example of her creative imagination."[169]
+
+Blyton's Malory Towers has been adapted into a musical of the same name by Emma Rice's theatre company. It was scheduled to do a UK spring tour in 2020 which has been postponed due to the COVID-19 pandemic.
+
+In 2019, Malory Towers was adapted as a 13 part TV series for the BBC. It is made partly in Toronto and partly in the UK in association with Canada's Family Channel. The series went to air in the UK from April 2020. [170]
+
+Seven Stories, the National Centre for Children's Books in Newcastle upon Tyne, holds the largest public collection of Blyton's papers and typescripts.[171] The Seven Stories collection contains a significant number of Blyton's typescripts, including the previously unpublished novel, Mr Tumpy's Caravan, as well as personal papers and diaries.[172] The purchase of the material in 2010 was made possible by special funding from the Heritage Lottery Fund, the MLA/V&A Purchase Grant Fund, and two private donations.[173]

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+Ennio Morricone, OMRI[1] (Italian: [ˈɛnnjo morriˈkoːne]; 10 November 1928 – 6 July 2020) was an Italian composer, orchestrator, conductor, and trumpet player who wrote music in a wide range of styles. Morricone composed over 400 scores for cinema and television, as well as over 100 classical works. His score to The Good, the Bad and the Ugly (1966) is considered one of the most influential soundtracks in history[2] and was inducted into the Grammy Hall of Fame.[3] His filmography includes over 70 award-winning films, all Sergio Leone's films since A Fistful of Dollars, all Giuseppe Tornatore's films since Cinema Paradiso, The Battle of Algiers, Dario Argento's Animal Trilogy, 1900, Exorcist II, Days of Heaven, several major films in French cinema, in particular the comedy trilogy La Cage aux Folles I, II, III and Le Professionnel, as well as The Thing, Once Upon a Time in America, The Mission, The Untouchables, Mission to Mars, Bugsy, Disclosure, In the Line of Fire, Bulworth, Ripley's Game and The Hateful Eight.[4] Morricone is widely regarded as one of the greatest and most influential film composers of all time.[5][6]
+
+After playing the trumpet in jazz bands in the 1940s, he became a studio arranger for RCA Victor and in 1955 started ghost writing for film and theatre. Throughout his career, he composed music for artists such as Paul Anka, Mina, Milva, Zucchero and Andrea Bocelli. From 1960 to 1975, Morricone gained international fame for composing music for Westerns and—with an estimated 10 million copies sold—Once Upon a Time in the West is one of the best-selling scores worldwide.[7] From 1966 to 1980, he was a main member of Il Gruppo, one of the first experimental composers collectives, and in 1969 he co-founded Forum Music Village, a prestigious recording studio. From the 1970s, Morricone excelled in Hollywood, composing for prolific American directors such as Don Siegel, Mike Nichols, Brian De Palma, Barry Levinson, Oliver Stone, Warren Beatty, John Carpenter and Quentin Tarantino. In 1977, he composed the official theme for the 1978 FIFA World Cup. He continued to compose music for European productions, such as Marco Polo, La piovra, Nostromo, Fateless, Karol and En mai, fais ce qu'il te plait. Morricone's music has been reused in television series, including The Simpsons and The Sopranos, and in many films, including Inglourious Basterds and Django Unchained. He also scored seven Westerns for Sergio Corbucci, Duccio Tessari's Ringo duology and Sergio Sollima's The Big Gundown and Face to Face. Morricone worked extensively for other film genres with directors such as Bernardo Bertolucci, Mauro Bolognini, Giuliano Montaldo, Roland Joffé, Roman Polanski, Henri Verneuil, Lucio Fulci, Umberto Lenzi and Pier Paolo Pasolini. His acclaimed soundtrack for The Mission (1986)[8] was certified gold in the United States. The album Yo-Yo Ma Plays Ennio Morricone stayed 105 weeks on the Billboard Top Classical Albums.[9]
+
+Morricone's best-known compositions include "The Ecstasy of Gold", "Se Telefonando", "Man with a Harmonica", "Here's to You", the UK No. 2 single "Chi Mai", "Gabriel's Oboe" and "E Più Ti Penso". In 1971, he received a "Targa d'Oro" for worldwide sales of 22 million,[10] and by 2016 Morricone had sold over 70 million records worldwide.[11] In 2007, he received the Academy Honorary Award "for his magnificent and multifaceted contributions to the art of film music." He was nominated for a further six Oscars, and in 2016, received his only competitive Academy Award for his score to Quentin Tarantino's film The Hateful Eight, at the time becoming the oldest person ever to win a competitive Oscar. His other achievements include three Grammy Awards, three Golden Globes, six BAFTAs, ten David di Donatello, eleven Nastro d'Argento, two European Film Awards, the Golden Lion Honorary Award and the Polar Music Prize in 2010. Morricone influenced many artists from film scoring to other styles and genres, including Hans Zimmer,[12] Danger Mouse,[13] Dire Straits,[14] Muse,[15] Metallica,[16] and Radiohead.[17]
+
+Morricone was born in Rome, the son of Libera Ridolfi and Mario Morricone, a musician. At the time of his birth Italy was under fascist rule.[18] His family came from Arpino, near Frosinone. Morricone had four siblings — Adriana, Aldo,[nb 1] Maria, and Franca — and lived in Trastevere in the centre of Rome. His father was a professional trumpet player who performed in light-music orchestras while his mother set up a small textile business.[19]
+
+Morricone's father first taught him to read music and to play several instruments. He entered the National Academy of Saint Cecilia to take trumpet lessons under the guidance of Umberto Semproni.[18] He formally entered the conservatory in 1940 at age 12, enrolling in a four-year harmony program which he completed within six months. He studied the trumpet, composition, and choral music under the direction of Goffredo Petrassi, to whom Morricone would later dedicate concert pieces.
+
+In 1941 Morricone was chosen among the students of the National Academy of Saint Cecilia to be a part of the Orchestra of the Opera directed by Carlo Zecchi on the occasion of a tour of the Veneto region.[20] He received his diploma in trumpet in 1946,[21] continuing to work in classical composition and arrangement.[18] Morricone received the Diploma in Instrumentation for Band Arrangement with a mark of 9/10 in 1952. His studies concluded at the Conservatory of Santa Cecilia in 1954 when he obtained a final 9.5/10 in his Diploma in Composition under Petrassi.[22]
+
+Morricone wrote his first compositions when he was six years old and was encouraged to develop his natural talents.[23] In 1946, he composed "Il Mattino" ("The Morning") for voice and piano on a text by Fukuko, first in a group of seven "youth" Lieder.[24]
+
+In the following years, he continued to write music for the theatre as well as classical music for voice and piano, such as "Imitazione", based on a text by Italian poet Giacomo Leopardi, "Intimità", based on a text by Olinto Dini, "Distacco I" and "Distacco II" with words by R. Gnoli, "Oboe Sommerso" for baritone and five instruments with words by poet Salvatore Quasimodo and "Verrà la Morte", for alto and piano, based on a text by novelist Cesare Pavese.[24]
+
+In 1953, Morricone was asked by Gorni Kramer and Lelio Luttazzi to write an arrangement for some medleys in an American style for a series of evening radio shows. The composer continued with the composition of other 'serious' classical pieces, thus demonstrating the flexibility and eclecticism which has always been an integral part of his character. Many orchestral and chamber compositions date, in fact, from the period between 1954 and 1959: Musica per archi e pianoforte (1954), Invenzione, Canone e Ricercare per piano; Sestetto per flauto, oboe, fagotto, violino, viola e violoncello (1955), Dodici Variazione per oboe, violoncello e piano; Trio per clarinetto, corno e violoncello; Variazione su un tema di Frescobaldi (1956); Quattro pezzi per chitarra (1957); Distanze per violino, violoncello e piano; Musica per undici violini, Tre Studi per flauto, clarinetto e fagotto (1958); and the Concerto per orchestra (1957), dedicated to his teacher Goffredo Petrassi.[24][25]
+
+Morricone soon gained popularity by writing his first background music for radio dramas and quickly moved into film.[26]
+
+Morricone's career as an arranger began in 1950, by arranging the piece Mamma Bianca (Narciso Parigi).[27] In occasion of the "Anno Santo" (Holy Year), he arranged a long group of popular songs of devotion for radio broadcasting.[28]
+
+In 1956, Morricone started to support his family by playing in a jazz band and arranging pop songs for the Italian broadcasting service RAI. He was hired by RAI in 1958, but quit his job on his first day at work when he was told that broadcasting of music composed by employees was forbidden by a company rule. Subsequently, Morricone became a top studio arranger at RCA Victor, working with Renato Rascel, Rita Pavone, Domenico Modugno and Mario Lanza.
+
+Throughout his career, Morricone composed songs for several national and international jazz and pop artists, including Gianni Morandi (Go Kart Twist, 1962), Alberto Lionello (La donna che vale, 1959), Edoardo Vianello (Ornella, 1960; Cicciona cha-cha, 1960; Faccio finta di dormire, 1961; T'ho conosciuta, 1963; ), Nora Orlandi (Arianna, 1960), Jimmy Fontana (Twist no. 9; Nicole, 1962), Rita Pavone (Come te non-ce nessuno and Pel di carota from 1962, arranged by Luis Bacalov), Catherine Spaak (Penso a te; Questi vent'anni miei, 1964), Luigi Tenco (Quello che conta; Tra tanta gente; 1962), Gino Paoli (Nel corso from 1963, written by Morricone with Paoli), Renato Rascel (Scirocco, 1964), Paul Anka (Ogni Volta), Amii Stewart, Rosy Armen (L'Amore Gira), Milva (Ridevi, Metti Una Sera A Cena), Françoise Hardy (Je changerais d'avis, 1966), Mireille Mathieu (Mon ami de toujours; Pas vu, pas pris, 1971; J'oublie la pluie et le soleil, 1974) and Demis Roussos (I Like The World, 1970).[29][30]
+
+In 1963, the composer co-wrote (with Roby Ferrante) the music for the composition "Ogni volta" ("Every Time"), a song that was performed by Paul Anka for the first time during the Festival di San Remo in 1964. This song was arranged and conducted by Morricone and sold over three million copies worldwide, including one million copies in Italy alone.[31]
+
+Another particular success was his composition, "Se telefonando."[32] Performed by Mina, it was a standout track of Studio Uno 66, the fifth-best-selling album of the year 1966 in Italy.[33] Morricone's sophisticated arrangement of "Se telefonando" was a combination of melodic trumpet lines, Hal Blaine–style drumming, a string set, a '60s Europop female choir, and intensive subsonic-sounding trombones. The Italian Hitparade No. 7 song had eight transitions of tonality building tension throughout the chorus. During the following decades, the song was recorded by several performers in Italy and abroad including covers by Françoise Hardy and Iva Zanicchi (1966), Delta V (2005), Vanessa and the O's (2007), and Neil Hannon (2008).[34] Françoise Hardy – Mon amie la rose site in the reader's poll conducted by the la Repubblica newspaper to celebrate Mina's 70th anniversary in 2010, 30,000 voters picked the track as the best song ever recorded by Mina.[35]
+
+In 1987, Morricone co-wrote 'It Couldn't Happen Here' with the Pet Shop Boys. Other compositions for international artists include: La metà di me and Immagina (1988) by Ruggero Raimondi, Libera l'amore (1989) performed by Zucchero, Love Affair (1994) by k.d. lang, Ha fatto un sogno (1997) by Antonello Venditti, Di Più (1997) by Tiziana Tosca Donati, Come un fiume tu (1998), Un Canto (1998) and Conradian (2006) by Andrea Bocelli, Ricordare (1998) and Salmo (2000) by Angelo Branduardi and My heart and I (2001) by Sting.[36]
+
+After graduating in 1954, Morricone started to write and arrange music as a ghost writer for films credited to already well-known composers, while also arranging for many light music orchestras of the RAI television network, working especially with Armando Trovajoli, Alessandro Cicognini and Carlo Savina. He occasionally adopted Anglicized pseudonyms, such as Dan Savio and Leo Nichols.[37]
+
+In 1959, Morricone was the conductor (and uncredited co-composer) for Mario Nascimbene's score to Morte Di Un Amico (Death of a Friend), an Italian drama directed by Franco Rossi. In the same year, he composed music for the theatre show Il Lieto Fine by Luciano Salce.
+
+The 1960s began on a positive note: 1961 marked in fact his real film debut with Luciano Salce's Il Federale (The Fascist). In an interview with American composer Fred Karlin, Morricone discussed his beginnings, stating, "My first films were light comedies or costume movies that required simple musical scores that were easily created, a genre that I never completely abandoned even when I went on to much more important films with major directors".[38]
+
+With Il Federale Morricone began a long-run collaboration with Luciano Salce. In 1962, Morricone composed the jazz-influenced score for Salce's comedy La voglia matta (Crazy Desire). That year Morricone arranged also Italian singer Edoardo Vianello's summer hit "Pinne, Fucile e Occhiali", a cha-cha song, peppered with added water effects, unusual instrumental sounds and unexpected stops and starts.[39]
+
+Morricone wrote works for the concert hall in a more avant-garde style.[40] Some of these have been recorded, such as Ut, a trumpet concerto dedicated to Mauro Maur.[41]
+
+From 1964 up to their eventual disbandment in 1980, Morricone was part of Gruppo di Improvvisazione Nuova Consonanza (G.I.N.C.), a group of composers who performed and recorded avant-garde free improvisations. The Rome-based avant-garde ensemble was dedicated to the development of improvisation and new music methods. The ensemble functioned as a laboratory of sorts, working with anti-musical systems and sound techniques in an attempt to redefine the new music ensemble and explore "New Consonance."[42]
+
+Known as "The Group" or "Il Gruppo," they released seven albums across the Deutsche Grammophon, RCA and Cramps labels: Gruppo di Improvvisazione Nuova Consonanza (1966), The Private Sea of Dreams (1967), Improvisationen (1968), The Feed-back (1970), Improvvisazioni a Formazioni Variate (1973), Nuova Consonanza (1975) and Musica su Schemi (1976). Perhaps the most famous of these is their album entitled The Feed-back, which combines free jazz and avant-garde classical music with funk; the album is frequently sampled by hip hop DJs and is considered to be one of the most collectable records in existence, often fetching over $1,000 at auction.[43]
+
+Morricone played a key role in The Group and was among the core members in its revolving line-up; in addition to serving as their trumpet player, he directed them on many occasions and they can be heard on a large number of his scores.[44] Held in high regard in avant-garde music circles, they are considered to be the first experimental composers collective, their only peers being the British improvisation collective AMM. Their influence can be heard in free improvising ensembles from the European movements including the Evan Parker Electro-Acoustic Ensemble, the Swiss electronic free improvisation group Voice Crack, John Zorn[45] and in the techniques of modern classical music and avant-garde jazz groups. The ensemble's groundbreaking work informed their work in composition. The ensemble also performed in varying capacities with Morricone, contributing to some of his '60s and '70s Italian soundtracks, including A Quiet Place in the Country (1969) and Cold Eyes of Fear (1971).[46]
+
+Morricone's earliest scores were Italian light comedy and costume pictures, where he learned to write simple, memorable themes. During the sixties and seventies he composed the scores for comedies such as Eighteen in the Sun (Diciottenni al sole, 1962), Il Successo (1963), Lina Wertmüller's I basilischi (The Basilisks/The Lizards, 1963),[38] Slalom (1965), Menage all'italiana (Menage Italian Style, 1965), How I Learned to Love Women (Come imparai ad amare le donne, 1966), Her Harem (L'harem, 1967), A Fine Pair (Ruba al prossimo tuo, 1968), L'Alibi (1969), This Kind of Love (Questa specie d'amore, 1972), Winged Devils (Forza "G", 1972) and Fiorina la vacca (1972).
+
+His best-known scores for comedies includes La Cage aux Folles (1978) and La Cage aux Folles II (1980), both directed by Édouard Molinaro, Il ladrone (The Good Thief, 1980), Georges Lautner's La Cage aux Folles 3: The Wedding (1985), Pedro Almodóvar's Tie Me Up! Tie Me Down! (1990) and Warren Beatty's Bulworth (1998). Morricone never ceased to arrange and write music for comedies. In 2007, he composed a lighthearted score for the Italian romantic comedy Tutte le Donne della mia Vita by Simona Izzo, the director who co-wrote the Morricone-scored religious mini-series Il Papa Buono.[47]
+
+Though his first films were undistinguished, Morricone's arrangement of an American folk song intrigued director and former schoolmate Sergio Leone. Before being associated with Leone, Morricone had already composed some music for less-known western movies such as Duello nel Texas (aka Gunfight at Red Sands) (1963). In 1962, Morricone met American folksinger Peter Tevis, with the two collaborating on a version of Woody Guthrie's Pastures of Plenty. Tevis is credited with singing the lyrics of Morricone's songs such as "A Gringo Like Me" (from Gunfight at Red Sands) and "Lonesome Billy" (from Bullets Don't Argue).[48] Tevis later recorded a vocal version of A Fistful of Dollars that was not used in the film.
+
+The turning point in Morricone's career took place in 1964, the year in which his third child, Andrea Morricone, who would also become a film composer, was born. Film director Sergio Leone hired Morricone, and together they created a distinctive score to accompany Leone's different version of the Western, A Fistful of Dollars (1964).[49]
+
+The Dollars Trilogy
+
+Because budget strictures limited Morricone's access to a full orchestra, he used gunshots, cracking whips, whistle, voices, jew's harp, trumpets, and the new Fender electric guitar, instead of orchestral arrangements of Western standards à la John Ford. Morricone used his special effects to punctuate and comically tweak the action—cluing in the audience to the taciturn man's ironic stance.[18]:69–77
+
+As memorable as Leone's close-ups, harsh violence, and black comedy, Morricone's work helped to expand the musical possibilities of film scoring. Morricone was initially billed on the film as Dan Savio. A Fistful of Dollars came out in Italy in 1964 and was released in America three years later, greatly popularising the so-called Spaghetti Western genre. For the American release, Sergio Leone and Ennio Morricone decided to adopt American-sounding names, so they called themselves respectively Bob Robertson and Dan Savio. Over the film's theatrical release, it grossed more than any other Italian film up to that point.[50] The film debuted in the United States in January 1967, where it grossed US$4.5 million for the year.[50] It eventually grossed $14.5 million in its American release,[50] against its budget of US$200,000.[51][52]
+
+With the score of A Fistful of Dollars, Morricone began his 20-year collaboration with his childhood friend Alessandro Alessandroni and his Cantori Moderni.[53] Alessandroni provided the whistling and the twanging guitar on the film scores, while his Cantori Moderni were a flexible troupe of modern singers. Morricone specifically exploited the solo soprano of the group, Edda Dell'Orso, at the height of her powers "an extraordinary voice at my disposal".[54]
+
+The composer subsequently scored Leone's other two Dollars Trilogy (or Man With No Name Trilogy) spaghetti westerns: For a Few Dollars More (1965) and The Good, the Bad and the Ugly (1966). All three films starred the American actor Clint Eastwood as The Man With No Name and depicted Leone's own intense vision of the mythical West. Morricone commented in 2007: "Some of the music was written before the film, which was unusual. Leone's films were made like that because he wanted the music to be an important part of it; he kept the scenes longer because he did not want the music to end." According to Morricone this explains "why the films are so slow."[55]
+
+Despite the small film budgets, the Dollars Trilogy was a box-office success. The available budget for The Good, the Bad and The Ugly was about US$1.2 million, but it became the most successful film of the Dollars Trilogy, grossing US$25.1 million in the United States and over 2,3 billion lire (1,2 million EUR) in Italy alone. Morricone's score became a major success and sold over three million copies worldwide, earning him over 200 million dollars. On 14 August 1968 the original score was certified by the RIAA with a golden record for the sale of 500,000 copies in the United States alone.[56]
+
+"The Ecstasy of Gold" became one of Morricone's best-known compositions. The opening scene of Jeff Tremaine's Jackass Number Two (2006), in which the cast is chased through a suburban neighbourhood by bulls, is accompanied by this piece. While punk rock band the Ramones used "The Ecstasy of Gold" as closing theme during their live performances, Metallica uses "The Ecstasy of Gold" as the introductory music for its concerts since 1983[57][58] This composition is also included on Metallica's live symphonic album S&M as well as the live album Live Shit: Binge & Purge. An instrumental metal cover by Metallica (with minimal vocals by lead singer James Hetfield) appeared on the 2007 Morricone tribute album We All Love Ennio Morricone. This metal version was nominated for a Grammy Award in the category of Best Rock Instrumental Performance. In 2009, the Grammy Award-winning hip-hop artist Coolio extensively sampled the theme for his song "Change".[59]
+
+Subsequent to the success of the Dollars trilogy, Morricone also composed the scores for Once Upon a Time in the West (1968) and Leone's last credited western film A Fistful of Dynamite (1971),[60] as well as the score for My Name Is Nobody (1973).[61]
+
+Morricone's score for Once Upon a Time in the West is one of the best-selling original instrumental scores in the world today, with up to 10 million copies sold, including one million copies[62][63] in France and over 800,000 copies[64] in the Netherlands.[65][66] One of the main themes from the score, "A Man with Harmonica" (L'uomo Dell'armonica), became known worldwide and sold over 1,260,000 copies in France.[67]
+
+The collaboration with Leone is considered one of the exemplary collaborations between a director and a composer. Morricone's last score for Leone was for his last film, the gangster drama Once Upon a Time in America (1984). Leone died on 30 April 1989 of a heart attack at the age of 60. Before his death in 1989, Leone was part-way through planning a film on the Siege of Leningrad, set during World War II. By 1989, Leone had been able to acquire US$100 million in financing from independent backers for the war epic. He had convinced Morricone to compose the film score. The project was cancelled when Leone died two days before he was to officially sign on for the film. In early 2003, Italian filmmaker Giuseppe Tornatore announced he would direct a film called Leningrad.[68] The film has yet to go into production and Morricone was cagey as to details on account of Tornatore's superstitious nature.[69]
+
+Two years after the start of his collaboration with Sergio Leone, Morricone also started to score music for another Spaghetti Western director, Sergio Corbucci. The composer wrote music for Corbucci's Navajo Joe (1966), The Hellbenders (1967), The Mercenary/The Professional Gun (1968), The Great Silence (1968), Compañeros (1970), Sonny and Jed (1972) and What Am I Doing in the Middle of the Revolution? (1972).[70][71]
+
+In addition, Morricone composed music for the western films by Sergio Sollima, The Big Gundown (with Lee Van Cleef, 1966), Face to Face (1967) and Run, Man, Run (1968), as well as the 1970 crime thriller Violent City (with Charles Bronson) and the poliziottesco film Revolver (1973).[70][72][73]
+
+Other relevant scores for less popular Spaghetti Westerns include Duello nel Texas (1963), Bullets Don't Argue (1964), A Pistol for Ringo (1965), The Return of Ringo (1965), Seven Guns for the MacGregors (1966), The Hills Run Red (1966), Giulio Petroni's Death Rides a Horse (1967) and Tepepa (1968), A Bullet for the General (1967), Guns for San Sebastian (with Charles Bronson and Anthony Quinn, 1968), A Sky Full of Stars for a Roof (1968), The Five Man Army (1969), Don Siegel's Two Mules for Sister Sara (1970), Life Is Tough, Eh Providence? (1972) and Buddy Goes West (1981).[18]:115–117
+
+With Leone's films, Ennio Morricone's name had been put firmly on the map. Most of Morricone's film scores of the 1960s were composed outside the Spaghetti Western genre, while still using Alessandroni's team. Their music included the themes for Il Malamondo (1964), Slalom (1965) and Listen, Let's Make Love (1967). In 1968, Morricone reduced his work outside the movie business and wrote scores for 20 films in the same year. The scores included psychedelic accompaniment for Mario Bava's superhero romp Danger: Diabolik (1968).[74]
+
+Morricone collaborated with Marco Bellocchio (Fists in the Pocket, 1965), Gillo Pontecorvo (The Battle of Algiers (1966) and Queimada! (1969) with Marlon Brando), Roberto Faenza (H2S, 1968), Giuliano Montaldo (Sacco e Vanzetti, 1971), Giuseppe Patroni Griffi ('Tis Pity She's a Whore, 1971), Mauro Bolognini (Drama of the Rich, 1974), Umberto Lenzi (Almost Human, 1974), Pier Paolo Pasolini (Salò, or the 120 Days of Sodom, 1975), Bernardo Bertolucci (Novecento, 1976) and Tinto Brass (The Key, 1983).[18]:115–116
+
+In 1970, Morricone wrote the score for Violent City. That same year, he received his first Nastro d'Argento for the music in Metti, una sera a cena (Giuseppe Patroni Griffi, 1969) and his second only a year later for Sacco e Vanzetti (Giuliano Montaldo, 1971), in which he collaborated with the legendary American folk singer and activist Joan Baez. His soundtrack for Sacco e Vanzetti contains another well-known composition by Morricone, the folk song "Here's to You", sung by Joan Baez. For the writing of the lyrics, Baez was inspired by a letter from Bartolomeo Vanzetti: "Father, yes, I am a prisoner / Fear not to relay my crime". The song became a hit in several countries, selling over 790,000 copies in France only.[75] The song was later included in movies such as The Life Aquatic with Steve Zissou.[76]
+
+In the beginning of the 1970s, Morricone achieved success with other singles, including A Fistful of Dynamite (1971) and God With Us (1974), having sold respectively 477,000 and 378,000 copies in France only.[77][78]
+
+Morricone's eclecticism found its way to films in the horror genre, such as the baroque thrillers of Dario Argento, from The Bird with the Crystal Plumage (1969), The Cat o' Nine Tails (1970) and Four Flies on Grey Velvet (1971) to The Stendhal Syndrome (1996) and The Phantom of the Opera (1998). His other horror scores include Nightmare Castle (1965), A Quiet Place in the Country (1968), The Antichrist (1974), Autopsy (1975) and Night Train Murders (1975).
+
+In addition, Morricone's music has also been featured in many popular and cult Italian giallo films, such as Senza sapere niente di lei (1969), Forbidden Photos of a Lady Above Suspicion (1970), A Lizard in a Woman's Skin (1971), Cold Eyes of Fear (1971), The Fifth Cord (1971), Short Night of Glass Dolls (1971), My Dear Killer (1972), What Have You Done to Solange? (1972), Black Belly of the Tarantula (1972), Who Saw Her Die? (1972) and Spasmo (1974).
+
+In 1977 Morricone scored John Boorman's Exorcist II: The Heretic and Alberto De Martino's apocalyptic horror film Holocaust 2000, starring Kirk Douglas. In 1982 he composed the score for John Carpenter's science fiction horror movie The Thing.[79] Morricone's main theme for the film was reflected in Marco Beltrami's film's score of prequel of the 1982 film, which was released in 2011.
+
+The Dollars Trilogy was not released in the United States until 1967 when United Artists, who had already enjoyed success distributing the British-produced James Bond films in the United States, decided to release Sergio Leone's Spaghetti Westerns. The American release gave Morricone an exposure in America and his film music became quite popular in the United States.[80]
+
+One of Morricone's first contributions to an American director concerned his music for the religious epic film The Bible: In the Beginning... by John Huston. According to Sergio Miceli's book Morricone, la musica, il cinema, Morricone wrote about 15 or 16 minutes of music, which were recorded for a screen test and conducted by Franco Ferrara. At first Morricone's teacher Goffredo Petrassi had been engaged to write the score for the great big budget epic, but Huston preferred another composer. RCA Records then proposed Morricone who was under contract with them, but a conflict between the film's producer Dino De Laurentiis and RCA occurred. The producer wanted to have the exclusive rights for the soundtrack, while RCA still had the monopoly on Morricone at that time and did not want to release the composer. Subsequently, Morricone's work was rejected because he did not get the permission by RCA to work for Dino De Laurentiis alone. The composer reused the parts of his unused score for The Bible: In the Beginning in such films as The Return of Ringo (1965) by Duccio Tessari and Alberto Negrin's The Secret of the Sahara (1987).
+
+Morricone never left Rome to compose his music and never learned to speak English. But given that the composer always worked in a wide field of composition genres, from "absolute music", which he always produced, to "applied music", working as orchestrator as well as conductor in the recording field, and then as a composer for theatre, radio and cinema, the impression arises that he never really cared that much about his standing in the eyes of Hollywood.[81]
+
+In 1970, Morricone composed the music for Don Siegel's Two Mules for Sister Sara, an American-Mexican western film starring Shirley MacLaine and Clint Eastwood. The same year the composer also delivered the title theme The Men from Shiloh for the American Western television series The Virginian.[82]
+
+In 1974-1975 Morricone wrote music for Spazio 1999, an Italian-produced compilation movie made to launch the Italian-British TV series Space: 1999,[83] while the original episodes featured music by Barry Gray. A soundtrack album was only released on CD in 2016[84] and on LP in 2017.[85] In 1975 he scored the George Kennedy revenge thriller The "Human" Factor, which was the final film of director Edward Dmytryk. Two years later he composed the score for the sequel to William Friedkin's 1973 film The Exorcist, directed by John Boorman: Exorcist II: The Heretic. The horror film was a major disappointment at the box office. The film grossed US$30,749,142 in the United States.[86]
+
+In 1978, the composer worked with Terrence Malick for Days of Heaven starring Richard Gere, for which he earned his first nomination at the Oscars for Best Original Score.[87]
+
+Despite the fact that Morricone had produced some of the most popular and widely imitated film music ever written throughout the 1960s and '70s, Days of Heaven earned him his first Oscar nomination for Best Original Score, with his score up against Jerry Goldsmith's The Boys from Brazil, Dave Grusin's Heaven Can Wait, Giorgio Moroder's Midnight Express (the eventual winner) and John Williams's Superman: The Movie at the Oscar ceremonies in 1979.[88]
+
+Association with Roland Joffé
+
+The Mission, directed by Joffé, was about a piece of history considerably more distant, as Spanish Jesuit missionaries see their work undone as a tribe of Paraguayan natives fall within a territorial dispute between the Spanish and Portuguese. At one point the score was one of the world's best-selling film scores, selling over 3 million copies worldwide.[89][90]
+
+Morricone finally received a second Oscar nomination for The Mission.[91] Morricone's original score lost out to Herbie Hancock's coolly arranged jazz on Bertrand Tavernier's Round Midnight. It was considered as a surprising win and a controversial one, given that much of the music in the film was pre-existing.[92] Morricone stated the following during a 2001 interview with The Guardian: "I definitely felt that I should have won for The Mission. Especially when you consider that the Oscar-winner that year was Round Midnight, which was not an original score. It had a very good arrangement by Herbie Hancock, but it used existing pieces. So there could be no comparison with The Mission. There was a theft!"[93] His score for The Mission was ranked at number 1 in a poll of the all-time greatest film scores. The top 10 list was compiled by 40 film composers such as Michael Giacchino and Carter Burwell.[92] The score is ranked 23rd on the AFI's list of 25 greatest film scores of all time.[94]
+
+On three occasions, Brian De Palma worked with Morricone: The Untouchables (1987), the 1989 war drama Casualties of War and the science fiction film Mission to Mars (2000).[79] Morricone's score for The Untouchables resulted in his third nomination for Academy Award for Best Original Score.[95]
+
+In a 2001 interview with The Guardian, Morricone stated that he had good experiences with De Palma: "De Palma is delicious! He respects music, he respects composers. For The Untouchables, everything I proposed to him was fine, but then he wanted a piece that I didn't like at all, and of course we didn't have an agreement on that. It was something I didn't want to write – a triumphal piece for the police. I think I wrote nine different pieces for this in total and I said, 'Please don't choose the seventh!' because it was the worst. And guess what he chose? The seventh one. But it really suits the movie."[93]
+
+Another American director, Barry Levinson, commissioned the composer on two occasions. First, for the crime-drama Bugsy, starring Warren Beatty, which received ten Oscar nominations,[96] winning two for Best Art Direction-Set Decoration (Dennis Gassner, Nancy Haigh) and Best Costume Design.[97]
+
+"He doesn't have a piano in his studio, I always thought that with composers, you sit at the piano, and you try to find the melody. There's no such thing with Morricone. He hears a melody, and he writes it down. He hears the orchestration completely done," said Levinson in an interview.[98]
+
+During his career in Hollywood, Morricone was approached for numerous other projects, including the Gregory Nava drama A Time of Destiny (1988),[99] Frantic by Polish-French director Roman Polanski (1988, starring Harrison Ford), Franco Zeffirelli's 1990 drama film Hamlet (starring Mel Gibson and Glenn Close), the neo-noir[100] crime film State of Grace by Phil Joanou (1990, starring Sean Penn and Ed Harris),[101] Rampage (1992) by William Friedkin,[102] and the romantic drama Love Affair (1994) by Warren Beatty.[103]
+
+In 2009, Tarantino originally wanted Morricone to compose the film score for Inglourious Basterds.[104][105] Morricone was unable to, because the film's sped-up production schedule conflicted with his scoring of Giuseppe Tornatore's Baarìa.[106] However, Tarantino did use eight tracks composed by Morricone in the film, with four of them included on the soundtrack. The tracks came originally from Morricone's scores for The Big Gundown (1966), Revolver (1973) and Allonsanfàn (1974).[107][108]
+
+In 2012, Morricone composed the song "Ancora Qui" with lyrics by Italian singer Elisa for Tarantino's Django Unchained, a track that appeared together with three existing music tracks composed by Morricone on the soundtrack. "Ancora Qui" was one of the contenders for an Academy Award nomination in the Best Original Song category, but eventually the song was not nominated.[109] On 4 January 2013 Morricone presented Tarantino with a Life Achievement Award at a special ceremony being cast as a continuation of the International Rome Film Festival.[110] In 2014, Morricone was misquoted, as claiming that he would "never work" with Tarantino again,[111]  and later agreed to write an original film score for Tarantino's The Hateful Eight, which won an Academy Award in 2016 in the Best Original Score category.[112]
+
+In 1988, Morricone started an ongoing and very successful collaboration with Italian director Giuseppe Tornatore. His first score for Tornatore was for the drama film Cinema Paradiso. The international version of the film won the Special Jury Prize at the 1989 Cannes Film Festival[113] and the 1989 Best Foreign Language Film Oscar. Morricone received a BAFTA award with his son Andrea, and a David di Donatello for his score. In 2002, the director's cut 173-minute version was released (known in the US as Cinema Paradiso: The New Version).[112]
+After the success of Cinema Paradiso, the composer wrote the music for all subsequent films by Tornatore: the drama film Everybody's Fine (Stanno Tutti Bene, 1990), A Pure Formality (1994) starring Gérard Depardieu and Roman Polanski, The Star Maker (1995), The Legend of 1900 (1998) starring Tim Roth, the 2000 romantic drama Malèna (which featured Monica Bellucci) and the psychological thriller mystery film La sconosciuta (2006). Morricone also composed the scores for Baarìa (2009), The Best Offer (2013) starring Geoffrey Rush, Jim Sturgess and Donald Sutherland and the romantic drama The Correspondence (2015)[112]
+
+The composer won several music awards for his scores to Tornatore's movies. So, Morricone received a fifth Academy Award nomination and a Golden Globe nomination for Malèna. For Legend of 1900, he won a Golden Globe Award for Best Original Score.[114]
+
+Morricone wrote the score for the Mafia television series La piovra seasons 2 to 10 from 1985 to 2001, including the themes "Droga e sangue" ("Drugs and Blood"), "La Morale", and "L'Immorale". Morricone worked as the conductor of seasons 3 to 5 of the series. He also worked as the music supervisor for the television project La bibbia ("The Bible").[115]
+In the late 1990s, he collaborated with his son Andrea on the Ultimo crime dramas, resulting in Ultimo (1998), Ultimo 2 – La sfida (1999), Ultimo 3 – L'infiltrato (2004) and Ultimo 4 – L'occhio del falco (2013).[116]
+In the 2000s, Morricone continued to compose music for successful television series such as Il Cuore nel Pozzo (2005), Karol: A Man Who Became Pope (2005), La provinciale (2006), Giovanni Falcone (2007), Pane e libertà (2009) and Come Un Delfino 1–2 (2011–2013).[117]
+
+Morricone provided the string arrangements on Morrissey's "Dear God Please Help Me" from the album Ringleader of the Tormentors in 2006.[118]
+
+In 2008, the composer recorded music for a Lancia commercial, featuring Richard Gere and directed by Harald Zwart (known for directing The Pink Panther 2).[119]
+
+In spring and summer 2010, Morricone worked with Hayley Westenra for a collaboration on her album Paradiso.[120] The album features new songs written by Morricone, as well as some of his best-known film compositions of the last 50 years.[121][122] Hayley recorded the album with Morricone's orchestra in Rome during the summer of 2010.[123][124][125]
+
+Since 1995, he composed the music for several advertising campaigns of Dolce & Gabbana. The commercials were directed by Giuseppe Tornatore.[126]
+
+In 2013, Morricone collaborated with Italian singer-songwriter Laura Pausini on a new version of her hit single "La solitudine" for her 20 years anniversary greatest hits album 20 – The Greatest Hits.[127]
+
+Morricone composed the music for The Best Offer (2013) by Giuseppe Tornatore.[128]
+
+He wrote the score for Christian Carion's En mai, fais ce qu'il te plait (2015) and the most recent movie by Tornatore: The Correspondence (2016), featuring Jeremy Irons and Olga Kurylenko.[129] In July 2015, Quentin Tarantino announced after the screening of footage of his movie The Hateful Eight at the San Diego Comic-Con International that Morricone would score the film, the first Western that Morricone scored since 1981.[130] The score was critically acclaimed and won several awards including the Golden Globe Award for Best Original Score and the Academy Award for Best Original Score.[131][132]
+
+Before receiving his diplomas in trumpet, composition and instrumentation from the conservatory, Morricone was already active as a trumpet player, often performing in an orchestra that specialised in music written for films. After completing his education at Saint Cecilia, the composer honed his orchestration skills as an arranger for Italian radio and television. In order to support himself, he moved to RCA in the early sixties and entered the front ranks of the Italian recording industry.[133] Since 1964, Morricone was also a founding member of the Rome-based avant-garde ensemble Gruppo di Improvvisazione Nuova Consonanza. During the existence of the group (until 1978), Morricone performed several times with the group as trumpet player.[134]
+
+To ready his music for live performance, he joined smaller pieces of music together into longer suites. Rather than single pieces, which would require the audience to applaud every few minutes, Morricone thought the best idea was to create a series of suites lasting from 15 to 20 minutes, which form a sort of symphony in various movements – alternating successful pieces with personal favourites. In concert, Morricone normally had 180 to 200 musicians and vocalists under his baton, performing multiple genre-crossing collections of music. Rock, symphonic and ethnic instruments share the stage.[135]
+
+On 20 September 1984 Morricone conducted the Orchestre national des Pays de la Loire at Cinésymphonie '84 ("Première nuit de la musique de film/First night of film music") in the French concert hall Salle Pleyel in Paris. He performed some of his best-known compositions such as Metti, una sera a cena, Novecento and The Good, the Bad and the Ugly.[136] Michel Legrand and Georges Delerue performed on the same evening.[137]
+
+On 15 October 1987 Morricone gave a concert in front of 12,000 people in the Sportpaleis in Antwerp, Belgium, with the Dutch Metropole Orchestra and the Italian operatic soprano Alide Maria Salvetta.[138] A live-album with a recording of this concert was released in the same year.[139]
+
+On 9 June 2000 Morricone went to the Flanders International Film Festival Ghent to conduct his music together with the National Orchestra of Belgium.[140] During the concert's first part, the screening of The Life and Death of King Richard III (1912) was accompanied with live music by Morricone. It was the very first time that the score was performed live in Europe. The second part of the evening consisted of an anthology of the composer's work. The event took place on the eve of Euro 2000, the European Football Championship in Belgium and the Netherlands.[141]
+
+Morricone performed over 250 concerts as of 2001.[142] The composer started a world tour in 2001, the latter part sponsored by Giorgio Armani, with the Orchestra Roma Sinfonietta, touring London (Barbican 2001; 75th birthday Concerto, Royal Albert Hall 2003), Paris, Verona, and Tokyo. Morricone performed his classic film scores at the Gasteig in Munich in 2004.[143]
+
+He made his North American concert debut on 3 February 2007 at Radio City Music Hall in New York City. The previous evening, Morricone had already presented at the United Nations a concert comprising some of his film themes, as well as the cantata Voci dal silenzio to welcome the new Secretary-General Ban Ki-Moon. A Los Angeles Times review bemoaned the poor acoustics and opined of Morricone: "His stick technique is adequate, but his charisma as a conductor is zero."[144]
+
+On 22 December 2012 Morricone conducted the 85-piece Belgian orchestra "Orkest der Lage Landen" and a 100-piece choir during a two-hour concert in the Sportpaleis in Antwerp.[145]
+
+In November 2013 Morricone began a world tour to coincide with the 50th anniversary of his film music career and performed in locations such as the Crocus City Hall in Moscow, Santiago, Chile, Berlin, Germany (O2 World, Germany), Budapest, Hungary, and Vienna (Stadhalle). Back in June 2014, Morricone had to cancel a US tour in New York (Barclays Center) and Los Angeles (Nokia Theatre LA Live) due to a back procedure on 20 February. Morricone postponed the rest of his world tour.[146]
+
+In November 2014 Morricone stated that he would resume his European tour starting from February 2015.[147]
+
+On 13 October 1956, Morricone married Maria Travia, whom he had met in 1950. Travia wrote lyrics to complement her husband's pieces. Her works include the Latin texts for The Mission. They had three sons and a daughter: Marco (1957), Alessandra (1961), the conductor and film composer Andrea (1964), and Giovanni Morricone (1966), a filmmaker, who lives in New York City.[148]
+
+Morricone lived in Italy his entire life and never desired to live in Hollywood. The New York Times Magazine listed him among hundreds of artists whose material was reportedly destroyed in the 2008 Universal fire.[149]
+
+Morricone described himself as a Christian leftist,[150] stating that he voted for the Christian Democracy (DC) for more than 40 years[151] and then, after its dissolution in 1994, he approached the centre-left coalition.[152]
+
+On 6 July 2020, Morricone died at the Università Campus Bio-Medico in Rome, aged 91, as a result of injuries sustained during a fall.[153][154]
+
+Ennio Morricone influenced many artists from other styles and genres, including Danger Mouse,[13] Dire Straits,[14] Muse,[15] Metallica,[16] Radiohead[17] and Hans Zimmer.[12]
+
+Morricone sold well over 70 million records worldwide during his career that spanned over seven decades,[168][169] including 6.5 million albums and singles in France,[170] over three million in the United States and more than two million albums in South Korea.[171] In 1971, the composer received his first golden record (disco d'oro) for the sale of 1,000,000 records in Italy[172][173] and a "Targa d'Oro" for the worldwide sales of 22 million.[10]
+
+Morricone received his first Academy Award nomination in 1979 for the score to Days of Heaven (Terrence Malick, 1978).[91] He received his second Oscar nomination for The Mission.[91] He also received Oscar nominations for his scores to The Untouchables (1987), Bugsy (1991), Malèna (2000), and The Hateful Eight (2016).[91] In February 2016, Morricone won his first competitive Academy Award for his score to The Hateful Eight.[174] Morricone and Alex North are the only composers to receive the Academy Honorary Award since its introduction in 1928.[175] He received the award in February 2007, "for his magnificent and multifaceted contributions to the art of film music."[176]
+In 2005, four film scores by Ennio Morricone were nominated by the American Film Institute for an honoured place in the AFI's Top 25 of Best American Film Scores of All Time.[177] His score for The Mission was ranked 23rd in the Top 25 list.[178]
+Morricone was nominated seven times for a Grammy Award. In 2009 The Recording Academy inducted his score for The Good, the Bad and the Ugly (1966) into the Grammy Hall of Fame.[179][180]
+In 2010 Ennio Morricone and Icelandic singer Björk have won the Polar Music Prize. The Polar Music Prize is Sweden's biggest music award and is typically shared by a pop artist and a classical musician. It was founded by Stig Anderson, manager of Swedish pop group ABBA, in 1989.

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+A secondary school describes an institution that provides secondary education and also usually includes the building where this takes place. Some secondary schools provide both lower secondary education (12 to 15 years of age) and upper secondary education (16 to 18 years of age) ie levels 2 and 3 of the ISCED scale, but these can also be provided in separate schools, as in the American middle and high school system. In the UK, elite public schools typically admit pupils between 13 and 18 years of age. UK state schools accommodate pupils between 11 and 18 years of age.
+
+Secondary schools follow on from primary schools and prepare for vocational or tertiary education. Attendance is usually compulsory for students until the age of 16. The organisations, buildings, and terminology are more or less unique in each country.[1][2]
+
+In the ISCED 2011 education scale levels 2 and 3 correspond to secondary education which are as follows:
+
+Within the English speaking world, there are three widely used systems to describe the age of the child. The first is the 'equivalent ages', then countries that base their education systems on the 'English model' use one of two methods to identify the year group, while countries that base their systems on the 'American K-12 model' refer to their year groups as 'grades'. The Irish model is structured similarly to the English model, but have significant differences in terms of labels. This terminology extends into research literature. Below is a convenient comparison [3]
+
+School building design does not happen in isolation. The building (or school campus) needs to accommodate:
+
+Each country will have  a different education system and priorities. [4] Schools need to accommodate students, staff, storage, mechanical and electrical systems, support staff, ancillary staff and administration. The number of rooms required can be determined from the predicted roll of the school and the area needed.
+
+According to standards used in the United Kingdom, a general classroom for 30 students needs to be 55 m2, or more generously 62 m2. A general art room for 30 students needs to be 83 m2, but 104 m2 for 3D textile work. A drama studio or a specialist science laboratory for 30 needs to be 90 m2. Examples are given on how this can be configured for a 1,200 place secondary (practical specialism).[5] and 1,850 place secondary school.[6]
+
+The building providing the education has to fulfil the needs of: The students, the teachers, the non-teaching support staff, the administrators and the community. It has to meet general government building guidelines, health requirements, minimal functional requirements for classrooms, toilets and showers, electricity and services, preparation and storage of textbooks and basic teaching aids. [7] An optimum secondary school will meet the minimum conditions and will have:
+
+Government accountants having read the advice then publish minimum guidelines on schools. These  enable environmental modelling and  establishing building costs. Future design plans are audited to ensure that these standards are met but not exceeded. Government ministries continue to press for the 'minimum' space and cost standards to be reduced.
+
+The UK government published this downwardly revised space formula in 2014. It said the floor area should be 1050m2 (+ 350m2 if there is a sixth form) + 6.3m2/pupil place for 11- to 16-year-olds + 7m2/pupil place for post-16s. The external finishes were to be downgraded to meet a build cost of £1113/m2. [8]
+
+A secondary school locally may be called  high school or senior high school. In some countries there are two phases to secondary education (ISCED 2) and (ISCED 3), here the junior high school, intermediate school, lower secondary school, or middle school occurs between the primary school (ISCED 1) and high school.
+

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+
+
+
+
+In meteorology, a cloud is an aerosol consisting of a visible mass of minute liquid droplets, frozen crystals, or other particles suspended in the atmosphere of a planetary body or similar space.[1] Water or various other chemicals may compose the droplets and crystals. On Earth, clouds are formed as a result of saturation of the air when it is cooled to its dew point, or when it gains sufficient moisture (usually in the form of water vapor) from an adjacent source to raise the dew point to the ambient temperature.
+
+They are seen in the Earth's homosphere, which includes the troposphere, stratosphere, and mesosphere. Nephology is the science of clouds, which is undertaken in the cloud physics branch of meteorology.  There are two methods of naming clouds in their respective layers of the homosphere, Latin and common.
+
+Genus types in the troposphere, the atmospheric layer closest to Earth's surface, have Latin names due to the universal adoption of Luke Howard's nomenclature that was formally proposed in 1802.  It became the basis of a modern international system that divides clouds into five physical forms which can be divided or classified further into altitude levels to derive the ten basic genera. The main representative cloud types for each of these forms are stratus, cirrus, stratocumulus, cumulus, and cumulonimbus.  Low-level stratiform and stratocumuliform genera do not have any altitude-related prefixes.  However mid-level variants of the same physical forms are given the prefix alto- while high-level types carry the prefix cirro-. The other main forms never have prefixes indicating altitude level. Cirriform clouds are always high-level while cumuliform and cumulonimbiform clouds are classified formally as low-level. The latter are also more informally characterized as multi-level or vertical as indicated by the cumulo- prefix.  Most of the ten genera derived by this method of classification can be subdivided into species and further subdivided into varieties. Very low stratiform clouds that extend down to the Earth's surface are given the common names fog and mist, but have no Latin names.
+
+In the stratosphere and mesosphere, clouds have common names for their main types. They may have the appearance of  stratiform veils or sheets, cirriform wisps, or stratocumuliform bands or ripples. They are seen infrequently, mostly in the polar regions of Earth. Clouds have been observed in the atmospheres of other planets and moons in the Solar System and beyond. However, due to their different temperature characteristics, they are often composed of other substances such as methane, ammonia, and sulfuric acid, as well as water.
+
+Tropospheric clouds can have a direct effect on climate change on Earth.  They may reflect incoming rays from the sun which can contribute to a cooling effect where and when these clouds occur, or trap longer wave radiation that reflects back up from the Earth's surface which can cause a warming effect.  The altitude, form, and thickness of the clouds are the main factors that affect the local heating or cooling of Earth and the atmosphere.  Clouds that form above the troposphere are too scarce and too thin to have any influence on climate change.
+
+The tabular overview that follows is very broad in scope. It draws from several methods of cloud classification, both formal and informal, used in different levels of the Earth's homosphere by a number of cited authorities. Despite some differences in methodologies and terminologies, the classification schemes seen in this article can be harmonized by using an informal cross-classification of physical forms and altitude levels to derive the 10 tropospheric genera, the fog and mist that forms at surface level, and several additional major types above the troposphere. The cumulus genus includes four species that indicate vertical size and structure which can affect both forms and levels. This table should not be seen as a strict or singular classification, but as an illustration of how various major cloud types are related to each other and defined through a full range of altitude levels from Earth's surface to the "edge of space".
+
+The origin of the term "cloud" can be found in the Old English words clud or clod, meaning a hill or a mass of rock. Around the beginning of the 13th century, the word came to be used as a metaphor for rain clouds, because of the similarity in appearance between a mass of rock and cumulus heap cloud. Over time, the metaphoric usage of the word supplanted the Old English weolcan, which had been the literal term for clouds in general.[2][3]
+
+Ancient cloud studies were not made in isolation, but were observed in combination with other weather elements and even other natural sciences. Around 340 BC, Greek philosopher Aristotle wrote Meteorologica, a work which represented the sum of knowledge of the time about natural science, including weather and climate. For the first time, precipitation and the clouds from which precipitation fell were called meteors, which originate from the Greek word meteoros, meaning 'high in the sky'. From that word came the modern term meteorology, the study of clouds and weather. Meteorologica was based on intuition and simple observation, but not on what is now considered the scientific method. Nevertheless, it was the first known work that attempted to treat a broad range of meteorological topics in a systematic way, especially the hydrological cycle.[4]
+
+After centuries of speculative theories about the formation and behavior of clouds, the first truly scientific studies were undertaken by Luke Howard in England and Jean-Baptiste Lamarck in France. Howard was a methodical observer with a strong grounding in the Latin language, and used his background to classify the various tropospheric cloud types during 1802. He believed that the changing cloud forms in the sky could unlock the key to weather forecasting. Lamarck had worked independently on cloud classification the same year and had come up with a different naming scheme that failed to make an impression even in his home country of France because it used unusual French names for cloud types. His system of nomenclature included 12 categories of clouds, with such names as (translated from French) hazy clouds, dappled clouds, and broom-like clouds. By contrast, Howard used universally accepted Latin, which caught on quickly after it was published in 1803.[5] As a sign of the popularity of the naming scheme, German dramatist and poet Johann Wolfgang von Goethe composed four poems about clouds, dedicating them to Howard. An elaboration of Howard's system was eventually formally adopted by the International Meteorological Conference in 1891.[5] This system covered only the tropospheric cloud types, but the discovery of clouds above the troposphere during the late 19th century eventually led to the creation of separate classification schemes using common names for these very high clouds, which were still broadly similar to some cloud forms identiified in the troposphhere.[6]
+
+Terrestrial clouds can be found throughout most of the homosphere, which includes the troposphere, stratosphere, and mesosphere. Within these layers of the atmosphere, air can become saturated as a result of being cooled to its dew point or by having moisture added from an adjacent source.[7] In the latter case, saturation occurs when the dew point is raised to the ambient air temperature.
+
+Adiabatic cooling occurs when one or more of three possible lifting agents – convective, cyclonic/frontal, or orographic – cause a parcel of air containing invisible water vapor to rise and cool to its dew point, the temperature at which the air becomes saturated. The main mechanism behind this process is adiabatic cooling.[8] As the air is cooled to its dew point and becomes saturated, water vapor normally condenses to form cloud drops. This condensation normally occurs on cloud condensation nuclei such as salt or dust particles that are small enough to be held aloft by normal circulation of the air.[9][10]
+
+One agent is the convective upward motion of air caused by daytime solar heating at surface level.[9] Airmass instability allows for the formation of cumuliform clouds that can produce showers if the air is sufficiently moist.[11] On moderately rare occasions, convective lift can be powerful enough to penetrate the tropopause and push the cloud top into the stratosphere.[12]
+
+Frontal and cyclonic lift occur when stable air is forced aloft at weather fronts and around centers of low pressure by a process called convergence.[13] Warm fronts associated with extratropical cyclones tend to generate mostly cirriform and stratiform clouds over a wide area unless the approaching warm airmass is unstable, in which case cumulus congestus or cumulonimbus clouds are usually embedded in the main precipitating cloud layer.[14] Cold fronts are usually faster moving and generate a narrower line of clouds, which are mostly stratocumuliform, cumuliform, or cumulonimbiform depending on the stability of the warm airmass just ahead of the front.[15]
+
+A third source of lift is wind circulation forcing air over a physical barrier such as a mountain (orographic lift).[9] If the air is generally stable, nothing more than lenticular cap clouds form. However, if the air becomes sufficiently moist and unstable, orographic showers or thunderstorms may appear.[16]
+
+Along with adiabatic cooling that requires a lifting agent, three major nonadiabatic mechanisms exist for lowering the temperature of the air to its dew point. Conductive, radiational, and evaporative cooling require no lifting mechanism and can cause condensation at surface level resulting in the formation of fog.[17][18][19]
+
+Several main sources of water vapor can be added to the air as a way of achieving saturation without any cooling process: water or moist ground,[20][21][22] precipitation or virga,[23] and transpiration from plants[24]
+
+Tropospheric classification is based on a hierarchy of categories with physical forms and altitude levels at the top.[25][26] These are cross-classified into a total of ten genus types, most of which can be divided into species and further subdivided into varieties which are at the bottom of the hierarchy.[27]
+
+Clouds in the troposphere assume five physical forms based on structure and process of formation. These forms are commonly used for the purpose of satellite analysis.[25] They are given below in approximate ascending order of instability or convective activity.[28]
+
+Nonconvective stratiform clouds appear in stable airmass conditions and, in general, have flat, sheet-like structures that can form at any altitude in the troposphere.[29] The stratiform group is divided by altitude range into the genera cirrostratus (high-level), altostratus (mid-level), stratus (low-level), and nimbostratus (multi-level).[26] Fog is commonly considered a surface-based cloud layer.[16] The fog may form at surface level in clear air or it may be the result of a very low stratus cloud subsiding to ground or sea level. Conversely, low stratiform clouds result when advection fog is lifted above surface level during breezy conditions.
+
+Cirriform clouds in the troposphere are of the genus cirrus and have the appearance of detached or semimerged filaments. They form at high tropospheric altitudes in air that is mostly stable with little or no convective activity, although denser patches may occasionally show buildups caused by limited high-level convection where the air is partly unstable.[30] Clouds resembling cirrus can be found above the troposphere but are classified separately using common names.
+
+Clouds of this structure have both cumuliform and stratiform characteristics in the form of rolls, ripples, or elements.[31] They generally form as a result of limited convection in an otherwise mostly stable airmass topped by an inversion layer.[32] If the inversion layer is absent or higher in the troposphere, increased airmass instability may cause the cloud layers to develop tops in the form of turrets consisting of embedded cumuliform buildups.[33] The stratocumuliform group is divided into cirrocumulus (high-level), altocumulus (mid-level), and stratocumulus (low-level).[31]
+
+Cumuliform clouds generally appear in isolated heaps or tufts.[34][35] They are the product of localized but generally free-convective lift where no inversion layers are in the troposphere to limit vertical growth. In general, small cumuliform clouds tend to indicate comparatively weak instability. Larger cumuliform types are a sign of greater atmospheric instability and convective activity.[36] Depending on their vertical size, clouds of the cumulus genus type may be low-level or multi-level with moderate to towering vertical extent.[26]
+
+The largest free-convective clouds comprise the genus cumulonimbus, which have towering vertical extent. They occur in highly unstable air[9] and often have fuzzy outlines at the upper parts of the clouds that sometimes include anvil tops.[31] These clouds are the product of very strong convection that can penetrate the lower stratosphere.
+
+Tropospheric clouds form in any of three levels (formerly called étages) based on altitude range above the Earth's surface. The grouping of clouds into levels is commonly done for the purposes of cloud atlases, surface weather observations,[26] and weather maps.[37] The base-height range for each level varies depending on the latitudinal geographical zone.[26] Each altitude level comprises two or three genus-types differentiated mainly by physical form.[38][31]
+
+The standard levels and genus-types are summarised below in approximate descending order of the altitude at which each is normally based.[39] Multi-level clouds with significant vertical extent are separately listed and summarized in approximate ascending order of instability or convective activity.[28]
+
+High clouds form at altitudes of 3,000 to 7,600 m (10,000 to 25,000 ft) in the polar regions, 5,000 to 12,200 m (16,500 to 40,000 ft) in the temperate regions, and 6,100 to 18,300 m (20,000 to 60,000 ft) in the tropics.[26] All cirriform clouds are classified as high, thus constitute a single genus cirrus (Ci). Stratocumuliform and stratiform clouds in the high altitude range carry the prefix cirro-, yielding the respective genus names cirrocumulus (Cc) and cirrostratus (Cs). When limited-resolution satellite images of high clouds are analysed without supporting data from direct human observations, distinguishing between individual forms or genus types becomes impossible, and they are then collectively identified as high-type (or informally as cirrus-type, though not all high clouds are of the cirrus form or genus).[40]
+
+Nonvertical clouds in the middle level are prefixed by alto-, yielding the genus names altocumulus (Ac) for stratocumuliform types and altostratus (As) for stratiform types. These clouds can form as low as 2,000 m (6,500 ft) above surface at any latitude, but may be based as high as 4,000 m (13,000 ft) near the poles, 7,000 m (23,000 ft) at midlatitudes, and 7,600 m (25,000 ft) in the tropics.[26] As with high clouds, the main genus types are easily identified by the human eye, but distinguishing between them using satellite photography is not possible. Without the support of human observations, these clouds are usually collectively identified as middle-type on satellite images.[40]
+
+Low clouds are found from near the surface up to 2,000 m (6,500 ft).[26] Genus types in this level either have no prefix or carry one that refers to a characteristic other than altitude. Clouds that form in the low level of the troposphere are generally of larger structure than those that form in the middle and high levels, so they can usually be identified by their forms and genus types using satellite photography alone.[40]
+
+
+
+These clouds have low- to mid-level bases that form anywhere from near the surface to about 2,400 m (8,000 ft) and tops that can extend into the mid-altitude range and sometimes higher in the case of nimbostratus.
+
+This is a diffuse, dark grey, multi-level stratiform layer with great horizontal extent and usually moderate to deep vertical development. It lacks towering structure and looks feebly illuminated from the inside.[58] Nimbostratus normally forms from mid-level altostratus, and develops at least moderate vertical extent[59][60] when the base subsides into the low level during precipitation that can reach moderate to heavy intensity. It  achieves even greater vertical development when it simultaneously grows upward into the high level due to large-scale frontal or cyclonic lift.[61] The nimbo- prefix refers to its ability to produce continuous rain or snow over a wide area, especially ahead of a warm front.[62] This thick cloud layer may be accompanied by embedded towering cumuliform or cumulonimbiform types.[60][63] Meteorologists affiliated with the World Meteorological Organization (WMO) officially classify nimbostratus as mid-level for synoptic purposes while informally characterizing it as multi-level.[26] Independent meteorologists and educators appear split between those who largely follow the WMO model[59][60] and those who classify nimbostratus as low-level, despite its considerable vertical extent and its usual initial formation in the middle altitude range.[64][65]
+
+These very large cumuliform and cumulonimbiform types have similar low- to mid-level cloud bases as the multi-level and moderate vertical types, and tops that nearly always extend into the high levels.  They are required to be identified by their standard names or abbreviations in all aviation observations (METARS) and forecasts (TAFS) to warn pilots of possible severe weather and turbulence.[66]
+
+Genus types are commonly divided into subtypes called species that indicate specific structural details which can vary according to the stability and windshear characteristics of the atmosphere at any given time and location. Despite this hierarchy, a particular species may be a subtype of more than one genus, especially if the genera are of the same physical form and are differentiated from each other mainly by altitude or level. There are a few species, each of which can be associated with genera of more than one physical form.[72] The species types are grouped below according to the physical forms and genera with which each is normally associated. The forms, genera, and species are listed in approximate ascending order of instability or convective activity.[28]
+
+Of the stratiform group, high-level cirrostratus comprises two species. Cirrostratus nebulosus has a rather diffuse appearance lacking in structural detail.[73] Cirrostratus fibratus is a species made of semi-merged filaments that are transitional to or from cirrus.[74] Mid-level altostratus and multi-level nimbostratus always have a flat or diffuse appearance and are therefore not subdivided into species. Low stratus is of the species nebulosus[73] except when broken up into ragged sheets of stratus fractus (see below).[59][72][75]
+
+Cirriform clouds have three non-convective species that can form in mostly stable airmass conditions. Cirrus fibratus comprise filaments that may be straight, wavy, or occasionally twisted by non-convective wind shear.[74] The species uncinus is similar but has upturned hooks at the ends. Cirrus spissatus appear as opaque patches that can show light grey shading.[72]
+
+Stratocumuliform genus-types (cirrocumulus, altocumulus, and stratocumulus) that appear in mostly stable air have two species each. The stratiformis species normally occur in extensive sheets or in smaller patches where there is only minimal convective activity.[76] Clouds of the lenticularis species tend to have lens-like shapes tapered at the ends. They are most commonly seen as orographic mountain-wave clouds, but can occur anywhere in the troposphere where there is strong wind shear combined with sufficient airmass stability to maintain a generally flat cloud structure. These two species can be found in the high, middle, or low levels of the troposphere depending on the stratocumuliform genus or genera present at any given time.[59][72][75]
+
+The species fractus shows variable instability because it can be a subdivision of genus-types of different physical forms that have different stability characteristics. This subtype can be in the form of ragged but mostly stable stratiform sheets (stratus fractus) or small ragged cumuliform heaps with somewhat greater instability (cumulus fractus).[72][75][77] When clouds of this species are associated with precipitating cloud systems of considerable vertical and sometimes horizontal extent, they are also classified as accessory clouds under the name pannus (see section on supplementary features).[78]
+
+These species are subdivisions of genus types that can occur in partly unstable air. The species castellanus appears when a mostly stable stratocumuliform or cirriform layer becomes disturbed by localized areas of airmass instability, usually in the morning or afternoon. This results in the formation of cumuliform buildups of limited convection arising from a common stratiform base.[79] Castellanus resembles the turrets of a castle when viewed from the side, and can be found with stratocumuliform genera at any tropospheric altitude level and with limited-convective patches of high-level cirrus.[80] Tufted clouds of the more detached floccus species are subdivisions of genus-types which may be cirriform or stratocumuliform in overall structure. They are sometimes seen with cirrus, cirrocumulus, altocumulus, and stratocumulus.[81]
+
+A newly recognized species of stratocumulus or altocumulus has been given the name volutus, a roll cloud that can occur ahead of a cumulonimbus formation.[82] There are some volutus clouds that form as a consequence of interactions with specific geographical features rather than with a parent cloud. Perhaps the strangest geographically specific cloud of this type is the Morning Glory, a rolling cylindrical cloud that appears unpredictably over the Gulf of Carpentaria in Northern Australia. Associated with a powerful "ripple" in the atmosphere, the cloud may be "surfed" in glider aircraft.[83]
+
+More general airmass instability in the troposphere tends to produce clouds of the more freely convective cumulus genus type, whose species are mainly indicators of degrees of atmospheric instability and resultant vertical development of the clouds. A cumulus cloud initially forms in the low level of the troposphere as a cloudlet of the species humilis that shows only slight vertical development. If the air becomes more unstable, the cloud tends to grow vertically into the species mediocris, then congestus, the tallest cumulus species[72] which is the same type that the International Civil Aviation Organization refers to as 'towering cumulus'.[66]
+
+With highly unstable atmospheric conditions, large cumulus may continue to grow into cumulonimbus calvus (essentially a very tall congestus cloud that produces thunder), then ultimately into the species capillatus when supercooled water droplets at the top of the cloud turn into ice crystals giving it a cirriform appearance.[72][75]
+
+Genus and species types are further subdivided into varieties whose names can appear after the species name to provide a fuller description of a cloud. Some cloud varieties are not restricted to a specific altitude level or form, and can therefore be common to more than one genus or species.[84]
+
+All cloud varieties fall into one of two main groups. One group identifies the opacities of particular low and mid-level cloud structures and comprises the varieties translucidus (thin translucent), perlucidus (thick opaque with translucent or very small clear breaks), and opacus (thick opaque). These varieties are always identifiable for cloud genera and species with variable opacity. All three are associated with the stratiformis species of altocumulus and stratocumulus. However, only two varieties are seen with altostratus and stratus nebulosus whose uniform structures prevent the formation of a perlucidus variety. Opacity-based varieties are not applied to high clouds because they are always translucent, or in the case of cirrus spissatus, always opaque.[84][85]
+
+A second group describes the occasional arrangements of cloud structures into particular patterns that are discernible by a surface-based observer (cloud fields usually being visible only from a significant altitude above the formations). These varieties are not always present with the genera and species with which they are otherwise associated, but only appear when atmospheric conditions favor their formation. Intortus and vertebratus varieties occur on occasion with cirrus fibratus. They are respectively filaments twisted into irregular shapes, and those that are arranged in fishbone patterns, usually by uneven wind currents that favor the formation of these varieties. The variety radiatus is associated with cloud rows of a particular type that appear to converge at the horizon. It is sometimes seen with the fibratus and uncinus species of cirrus, the stratiformis species of altocumulus and stratocumulus, the mediocris and sometimes humilis species of cumulus,[87][88] and with the genus altostratus.[89]
+
+Another variety, duplicatus (closely spaced layers of the same type, one above the other), is sometimes found with cirrus of both the fibratus and uncinus species, and with altocumulus and stratocumulus of the species stratiformis and lenticularis. The variety undulatus (having a wavy undulating base) can occur with any clouds of the species stratiformis or lenticularis, and with altostratus. It is only rarely observed with stratus nebulosus. The variety lacunosus is caused by localized downdrafts that create circular holes in the form of a honeycomb or net. It is occasionally seen with cirrocumulus and altocumulus of the species stratiformis, castellanus, and floccus, and with stratocumulus of the species stratiformis and castellanus.[84][85]
+
+It is possible for some species to show combined varieties at one time, especially if one variety is opacity-based and the other is pattern-based. An example of this would be a layer of altocumulus stratiformis arranged in seemingly converging rows separated by small breaks. The full technical name of a cloud in this configuration would be altocumulus stratiformis radiatus perlucidus, which would identify respectively its genus, species, and two combined varieties.[75][84][85]
+
+Supplementary features and accessory clouds are not further subdivisions of cloud types below the species and variety level. Rather, they are either hydrometeors or special cloud types with their own Latin names that form in association with certain cloud genera, species, and varieties.[75][85] Supplementary features, whether in the form of clouds or precipitation, are directly attached to the main genus-cloud. Accessory clouds, by contrast, are generally detached from the main cloud.[90]
+
+One group of supplementary features are not actual cloud formations, but precipitation that falls when water droplets or ice crystals that make up visible clouds have grown too heavy to remain aloft. Virga is a feature seen with clouds producing precipitation that evaporates before reaching the ground, these being of the genera cirrocumulus, altocumulus, altostratus, nimbostratus, stratocumulus, cumulus, and cumulonimbus.[90]
+
+When the precipitation reaches the ground without completely evaporating, it is designated as the feature praecipitatio.[91] This normally occurs with altostratus opacus, which can produce widespread but usually light precipitation, and with thicker clouds that show significant vertical development. Of the latter, upward-growing cumulus mediocris produces only isolated light showers, while downward growing nimbostratus is capable of heavier, more extensive precipitation. Towering vertical clouds have the greatest ability to produce intense precipitation events, but these tend to be localized unless organized along fast-moving cold fronts. Showers of moderate to heavy intensity can fall from cumulus congestus clouds. Cumulonimbus, the largest of all cloud genera, has the capacity to produce very heavy showers. Low stratus clouds usually produce only light precipitation, but this always occurs as the feature praecipitatio due to the fact this cloud genus lies too close to the ground to allow for the formation of virga.[75][85][90]
+
+Incus is the most type-specific supplementary feature, seen only with cumulonimbus of the species capillatus. A cumulonimbus incus cloud top is one that has spread out into a clear anvil shape as a result of rising air currents hitting the stability layer at the tropopause where the air no longer continues to get colder with increasing altitude.[92]
+
+The mamma feature forms on the bases of clouds as downward-facing bubble-like protuberances caused by localized downdrafts within the cloud. It is also sometimes called mammatus, an earlier version of the term used before a standardization of Latin nomenclature brought about by the World Meteorological Organization during the 20th century. The best-known is cumulonimbus with mammatus, but the mamma feature is also seen occasionally with cirrus, cirrocumulus, altocumulus, altostratus, and stratocumulus.[90]
+
+A tuba feature is a cloud column that may hang from the bottom of a cumulus or cumulonimbus. A newly formed or poorly organized column might be comparatively benign, but can quickly intensify into a funnel cloud or tornado.[90][93][94]
+
+An arcus feature is a roll cloud with ragged edges attached to the lower front part of cumulus congestus or cumulonimbus that forms along the leading edge of a squall line or thunderstorm outflow.[95] A large arcus formation can have the appearance of a dark menacing arch.[90]
+
+Several new supplementary features have been formally recognized by the World Meteorological Organization (WMO). The feature fluctus can form under conditions of strong atmospheric wind shear when a stratocumulus, altocumulus, or cirrus cloud breaks into regularly spaced crests. This variant is sometimes known informally as a Kelvin–Helmholtz (wave) cloud. This phenomenon has also been observed in cloud formations over other planets and even in the sun's atmosphere.[96] Another highly disturbed but more chaotic wave-like cloud feature associated with stratocumulus or altocumulus cloud has been given the Latin name asperitas. The supplementary feature cavum is a circular fall-streak hole that occasionally forms in a thin layer of supercooled altocumulus or cirrocumulus. Fall streaks consisting of virga or wisps of cirrus are usually seen beneath the hole as ice crystals fall out to a lower altitude. This type of hole is usually larger than typical lacunosus holes. A murus feature is a cumulonimbus wall cloud with a lowering, rotating cloud base than can lead to the development of tornadoes. A cauda feature is a tail cloud that extends horizontally away from the murus cloud and is the result of air feeding into the storm.[82]
+
+Supplementary cloud formations detached from the main cloud are known as accessory clouds.[75][85][90] The heavier precipitating clouds, nimbostratus, towering cumulus (cumulus congestus), and cumulonimbus typically see the formation in precipitation of the pannus feature, low ragged clouds of the genera and species cumulus fractus or stratus fractus.[78]
+
+A group of accessory clouds comprise formations that are associated mainly with upward-growing cumuliform and cumulonimbiform clouds of free convection. Pileus is a cap cloud that can form over a cumulonimbus or large cumulus cloud,[97] whereas a velum feature is a thin horizontal sheet that sometimes forms like an apron around the middle or in front of the parent cloud.[90] An accessory cloud recently officially recognized the World meteorological Organization is the flumen, also known more informally as the beaver's tail. It is formed by the warm, humid inflow of a super-cell thunderstorm, and can be mistaken for a tornado. Although the flumen can indicate a tornado risk, it is similar in appearance to pannus or scud clouds and does not rotate.[82]
+
+Clouds initially form in clear air or become clouds when fog rises above surface level. The genus of a newly formed cloud is determined mainly by air mass characteristics such as stability and moisture content. If these characteristics change over time, the genus tends to change accordingly. When this happens, the original genus is called a mother cloud. If the mother cloud retains much of its original form after the appearance of the new genus, it is termed a genitus cloud. One example of this is stratocumulus cumulogenitus, a stratocumulus cloud formed by the partial spreading of a cumulus type when there is a loss of convective lift. If the mother cloud undergoes a complete change in genus, it is considered to be a mutatus cloud.[98]
+
+The genitus and mutatus categories have been expanded to include certain types that do not originate from pre-existing clouds. The term flammagenitus (Latin for 'fire-made') applies to cumulus congestus or cumulonimbus that are formed by large scale fires or volcanic eruptions. Smaller low-level "pyrocumulus" or "fumulus" clouds formed by contained industrial activity are now classified as cumulus homogenitus (Latin for 'man-made'). Contrails formed from the exhaust of aircraft flying in the upper level of the troposphere can persist and spread into formations resembling cirrus which are designated cirrus homogenitus.  If a cirrus homogenitus cloud changes fully to any of the high-level genera, they are termed cirrus, cirrostratus, or cirrocumulus homomutatus. Stratus cataractagenitus (Latin for 'cataract-made') are generated by the spray from waterfalls. Silvagenitus (Latin for 'forest-made') is a stratus cloud that forms as water vapor is added to the air above a forest canopy.[98]
+
+Stratocumulus clouds can be organized into "fields" that take on certain specially classified shapes and characteristics. In general, these fields are more discernible from high altitudes than from ground level. They can often be found in the following forms:
+
+These patterns are formed from a phenomenon known as a Kármán vortex which is named after the engineer and fluid dynamicist Theodore von Kármán,.[101] Wind driven clouds can form into parallel rows that follow the wind direction. When the wind and clouds encounter high elevation land features such as a vertically prominent islands, they can form eddies around the high land masses that give the clouds a twisted appearance.[102]
+
+Although the local distribution of clouds can be significantly influenced by topography, the global prevalence of cloud cover in the troposphere tends to vary more by latitude. It is most prevalent in and along low pressure zones of surface tropospheric convergence which encircle the Earth close to the equator and near the 50th parallels of latitude in the northern and southern hemispheres.[105] The adiabatic cooling processes that lead to the creation of clouds by way of lifting agents are all associated with convergence; a process that involves the horizontal inflow and accumulation of air at a given location, as well as the rate at which this happens.[106] Near the equator, increased cloudiness is due to the presence of the low-pressure Intertropical Convergence Zone (ITCZ) where very warm and unstable air promotes mostly cumuliform and cumulonimbiform clouds.[107] Clouds of virtually any type can form along the mid-latitude convergence zones depending on the stability and moisture content of the air. These extratropical convergence zones are occupied by the polar fronts where air masses of polar origin meet and clash with those of tropical or subtropical origin.[108] This leads to the formation of weather-making extratropical cyclones composed of cloud systems that may be stable or unstable to varying degrees according to the stability characteristics of the various airmasses that are in conflict.[109]
+
+Divergence is the opposite of convergence. In the Earth's troposphere, it involves the horizontal outflow of air from the upper part of a rising column of air, or from the lower part of a subsiding column often associated with an area or ridge of high pressure.[106] Cloudiness tends to be least prevalent near the poles and in the subtropics close to the 30th parallels, north and south. The latter are sometimes referred to as the horse latitudes. The presence of a large-scale high-pressure subtropical ridge on each side of the equator reduces cloudiness at these low latitudes.[110] Similar patterns also occur at higher latitudes in both hemispheres.[111]
+
+The luminance or brightness of a cloud is determined by how light is reflected, scattered, and transmitted by the cloud's particles. Its brightness may also be affected by the presence of haze or photometeors such as halos and rainbows.[112] In the troposphere, dense, deep clouds exhibit a high reflectance (70% to 95%) throughout the visible spectrum. Tiny particles of water are densely packed and sunlight cannot penetrate far into the cloud before it is reflected out, giving a cloud its characteristic white color, especially when viewed from the top.[113] Cloud droplets tend to scatter light efficiently, so that the intensity of the solar radiation decreases with depth into the gases. As a result, the cloud base can vary from a very light to very-dark-grey depending on the cloud's thickness and how much light is being reflected or transmitted back to the observer. High thin tropospheric clouds reflect less light because of the comparatively low concentration of constituent ice crystals or supercooled water droplets which results in a slightly off-white appearance. However, a thick dense ice-crystal cloud appears brilliant white with pronounced grey shading because of its greater reflectivity.[112]
+
+As a tropospheric cloud matures, the dense water droplets may combine to produce larger droplets. If the droplets become too large and heavy to be kept aloft by the air circulation, they will fall from the cloud as rain. By this process of accumulation, the space between droplets becomes increasingly larger, permitting light to penetrate farther into the cloud. If the cloud is sufficiently large and the droplets within are spaced far enough apart, a percentage of the light that enters the cloud is not reflected back out but is absorbed giving the cloud a darker look. A simple example of this is one's being able to see farther in heavy rain than in heavy fog. This process of reflection/absorption is what causes the range of cloud color from white to black.[114]
+
+Striking cloud colorations can be seen at any altitude, with the color of a cloud usually being the same as the incident light.[115] During daytime when the sun is relatively high in the sky, tropospheric clouds generally appear bright white on top with varying shades of grey underneath. Thin clouds may look white or appear to have acquired the color of their environment or background. Red, orange, and pink clouds occur almost entirely at sunrise/sunset and are the result of the scattering of sunlight by the atmosphere. When the sun is just below the horizon, low-level clouds are gray, middle clouds appear rose-colored, and high clouds are white or off-white. Clouds at night are black or dark grey in a moonless sky, or whitish when illuminated by the moon. They may also reflect the colors of large fires, city lights, or auroras that might be present.[115]
+
+A cumulonimbus cloud that appears to have a greenish or bluish tint is a sign that it contains extremely high amounts of water; hail or rain which scatter light in a way that gives the cloud a blue color. A green colorization occurs mostly late in the day when the sun is comparatively low in the sky and the incident sunlight has a reddish tinge that appears green when illuminating a very tall bluish cloud. Supercell type storms are more likely to be characterized by this but any storm can appear this way. Coloration such as this does not directly indicate that it is a severe thunderstorm, it only confirms its potential. Since a green/blue tint signifies copious amounts of water, a strong updraft to support it, high winds from the storm raining out, and wet hail; all elements that improve the chance for it to become severe, can all be inferred from this. In addition, the stronger the updraft is, the more likely the storm is to undergo tornadogenesis and to produce large hail and high winds.[116]
+
+Yellowish clouds may be seen in the troposphere in the late spring through early fall months during forest fire season. The yellow color is due to the presence of pollutants in the smoke. Yellowish clouds are caused by the presence of nitrogen dioxide and are sometimes seen in urban areas with high air pollution levels.[117]
+
+Stratocumulus stratiformis and small castellanus made orange by the sun rising
+
+An occurrence of cloud iridescence with altocumulus volutus and cirrocumulus stratiformis
+
+Sunset reflecting shades of pink onto grey stratocumulus stratiformis translucidus (becoming perlucidus in the background)
+
+Stratocumulus stratiformis perlucidus before sunset. Bangalore, India.
+
+Late-summer rainstorm in Denmark. Nearly black color of base indicates main cloud in foreground probably cumulonimbus.
+
+Particles in the atmosphere and the sun's angle enhance colors of stratocumulus cumulogenitus at evening twilight
+
+Tropospheric clouds exert numerous influences on Earth's troposphere and climate. First and foremost, they are the source of precipitation, thereby greatly influencing the distribution and amount of precipitation. Because of their differential buoyancy relative to surrounding cloud-free air, clouds can be associated with vertical motions of the air that may be convective, frontal, or cyclonic. The motion is upward if the clouds are less dense because condensation of water vapor releases heat, warming the air and thereby decreasing its density. This can lead to downward motion because lifting of the air results in cooling that increases its density. All of these effects are subtly dependent on the vertical temperature and moisture structure of the atmosphere and result in major redistribution of heat that affect the Earth's climate.[118]
+
+The complexity and diversity of clouds in the troposphere is a major reason for difficulty in quantifying the effects of clouds on climate and climate change. On the one hand, white cloud tops promote cooling of Earth's surface by reflecting shortwave radiation (visible and near infrared) from the sun, diminishing the amount of solar radiation that is absorbed at the surface, enhancing the Earth's albedo. Most of the sunlight that reaches the ground is absorbed, warming the surface, which emits radiation upward at longer, infrared, wavelengths. At these wavelengths, however, water in the clouds acts as an efficient absorber. The water reacts by radiating, also in the infrared, both upward and downward, and the downward longwave radiation results in increased warming at the surface. This is analogous to the greenhouse effect of greenhouse gases and water vapor.[118]
+
+High-level genus-types particularly show this duality with both short-wave albedo cooling and long-wave greenhouse warming effects. On the whole, ice-crystal clouds in the upper troposphere (cirrus) tend to favor net warming.[119][120] However, the cooling effect is dominant with mid-level and low clouds, especially when they form in extensive sheets.[119] Measurements by NASA indicate that on the whole, the effects of low and mid-level clouds that tend to promote cooling outweigh the warming effects of high layers and the variable outcomes associated with vertically developed clouds.[119]
+
+As difficult as it is to evaluate the influences of current clouds on current climate, it is even more problematic to predict changes in cloud patterns and properties in a future, warmer climate, and the resultant cloud influences on future climate. In a warmer climate more water would enter the atmosphere by evaporation at the surface; as clouds are formed from water vapor, cloudiness would be expected to increase. But in a warmer climate, higher temperatures would tend to evaporate clouds.[121] Both of these statements are considered accurate, and both phenomena, known as cloud feedbacks, are found in climate model calculations. Broadly speaking, if clouds, especially low clouds, increase in a warmer climate, the resultant cooling effect leads to a negative feedback in climate response to increased greenhouse gases. But if low clouds decrease, or if high clouds increase, the feedback is positive. Differing amounts of these feedbacks are the principal reason for differences in climate sensitivities of current global climate models. As a consequence, much research has focused on the response of low and vertical clouds to a changing climate. Leading global models produce quite different results, however, with some showing increasing low clouds and others showing decreases.[122][123] For these reasons the role of tropospheric clouds in regulating weather and climate remains a leading source of uncertainty in global warming projections.[124][125]
+
+Polar stratospheric clouds (PSC's) form in the lowest part of the stratosphere during the winter, at the altitude and during the season that produces the coldest temperatures and therefore the best chances of triggering condensation caused by adiabatic cooling.   Moisture is scarce in the stratosphere, so nacreous and non-nacreous cloud at this altitude range is restricted to polar regions in the winter where the air is coldest.[6]
+
+PSC's show some variation in structure according to their chemical makeup and atmospheric conditions, but are limited to a single very high range of altitude of about 15,000–25,000 m (49,200–82,000 ft), so they are not classified into altitude levels, genus types, species, or varieties. There is no Latin nomenclature in the manner of tropospheric clouds, but rather descriptive names using common English.[6]
+
+Supercooled nitric acid and water PSC's, sometimes known as type 1, typically have a stratiform appearance resembling cirrostratus or haze, but because they are not frozen into crystals, do not show the pastel colours of the nacreous types.  This type of PSC has been identified as a cause of ozone depletion in the stratosphere.[126] The frozen nacreous types are typically very thin with mother-of-pearl colorations and an undulating cirriform or lenticular (stratocumuliform) appearance.  These are sometimes known as type 2.[127][128]
+
+Polar mesospheric clouds form at an extreme-level altitude range of about 80 to 85 km (50 to 53 mi). They are given the Latin name noctilucent because of their illumination well after sunset and before sunrise. They typically have a bluish or silvery white coloration that can resemble brightly illuminated cirrus. Noctilucent clouds may occasionally take on more of a red or orange hue.[6] They are not common or widespread enough to have a significant effect on climate.[129] However, an increasing frequency of occurrence of noctilucent clouds since the 19th century may be the result of climate change.[130]
+
+Noctilucent clouds are the highest in the atmosphere and form near the top of the mesosphere at about ten times the altitude of tropospheric high clouds.[131] From ground level, they can occasionally be seen illuminated by the sun during deep twilight. Ongoing research indicates that convective lift in the mesosphere is strong enough during the polar summer to cause adiabatic cooling of small amount of water vapour to the point of saturation. This tends to produce the coldest temperatures in the entire atmosphere just below the mesopause. These conditions result in the best environment for the formation of polar mesospheric clouds.[129] There is also evidence that smoke particles from burnt-up meteors provide much of the condensation nuclei required for the formation of noctilucent cloud.[132]
+
+Noctilucent clouds have four major types based on physical structure and appearance. Type I veils are very tenuous and lack well-defined structure, somewhat like cirrostratus or poorly defined cirrus.[133] Type II bands are long streaks that often occur in groups arranged roughly parallel to each other. They are usually more widely spaced than the bands or elements seen with cirrocumulus clouds.[134] Type III billows are arrangements of closely spaced, roughly parallel short streaks that mostly resemble cirrus.[135] Type IV whirls are partial or, more rarely, complete rings of cloud with dark centres.[136]
+
+Distribution in the mesosphere is similar to the stratosphere except at much higher altitudes. Because of the need for maximum cooling of the water vapor to produce noctilucent clouds, their distribution tends to be restricted to polar regions of Earth. A major seasonal difference is that convective lift from below the mesosphere pushes very scarce water vapor to higher colder altitudes required for cloud formation during the respective summer seasons in the northern and southern hemispheres. Sightings are rare more than 45 degrees south of the north pole or north of the south pole.[6]
+
+Cloud cover has been seen on most other planets in the Solar System. Venus's thick clouds are composed of sulfur dioxide (due to volcanic activity) and appear to be almost entirely stratiform.[137] They are arranged in three main layers at altitudes of 45 to 65 km that obscure the planet's surface and can produce virga. No embedded cumuliform types have been identified, but broken stratocumuliform wave formations are sometimes seen in the top layer that reveal more continuous layer clouds underneath.[138] On Mars, noctilucent, cirrus, cirrocumulus and stratocumulus composed of water-ice have been detected mostly near the poles.[139][140] Water-ice fogs have also been detected on Mars.[141]
+
+Both Jupiter and Saturn have an outer cirriform cloud deck composed of ammonia,[142][143] an intermediate stratiform haze-cloud layer made of ammonium hydrosulfide, and an inner deck of cumulus water clouds.[144][145] Embedded cumulonimbus are known to exist near the Great Red Spot on Jupiter.[146][147] The same category-types can be found covering Uranus, and Neptune, but are all composed of methane.[148][149][150][151][152][153] Saturn's moon Titan has cirrus clouds believed to be composed largely of methane.[154][155] The Cassini–Huygens Saturn mission uncovered evidence of polar stratospheric clouds[156] and a methane cycle on Titan, including lakes near the poles and fluvial channels on the surface of the moon.[157]
+
+Some planets outside the Solar System are known to have atmospheric clouds. In October 2013, the detection of high altitude optically thick clouds in the atmosphere of exoplanet Kepler-7b was announced,[158][159] and, in December 2013, in the atmospheres of GJ 436 b and GJ 1214 b.[160][161][162][163]
+
+Clouds play an important role in various cultures and religious traditions. The ancient Akkadians believed that the clouds were the breasts of the sky goddess Antu[165] and that rain was milk from her breasts.[165] In Exodus 13:21–22, Yahweh is described as guiding the Israelites through the desert in the form of a "pillar of cloud" by day and a "pillar of fire" by night.[164]
+
+In the ancient Greek comedy The Clouds, written by Aristophanes and first performed at the City Dionysia in 423 BC, the philosopher Socrates declares that the Clouds are the only true deities[166] and tells the main character Strepsiades not to worship any deities other than the Clouds, but to pay homage to them alone.[166] In the play, the Clouds change shape to reveal the true nature of whoever is looking at them,[167][166][168] turning into centaurs at the sight of a long-haired politician, wolves at the sight of the embezzler Simon, deer at the sight of the coward Cleonymus, and mortal women at the sight of the effeminate informer Cleisthenes.[167][168][166] They are hailed the source of inspiration to comic poets and philosophers;[166] they are masters of rhetoric, regarding eloquence and sophistry alike as their "friends".[166]
+
+In China, clouds are symbols of luck and happiness.[169] Overlapping clouds are thought to imply eternal happiness[169] and clouds of different colors are said to indicate "multiplied blessings".[169]

en/1770.html.txt

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+
+
+An integer (from the Latin integer meaning "whole")[a] is colloquially defined as a number that can be written without a fractional component. For example, 21, 4, 0, and −2048 are integers, while 9.75, 5+1/2, and √2 are not.
+
+The set of integers consists of zero (0), the positive natural numbers (1, 2, 3, ...), also called whole numbers or counting numbers,[2][3] and their additive inverses (the negative integers, i.e., −1, −2, −3, ...). The set of integers is often denoted by a boldface letter ‘Z’ ("Z") or blackboard bold 
+
+
+
+
+Z
+
+
+
+{\displaystyle \mathbb {Z} }
+
+ (Unicode U+2124 ℤ) standing for the German word Zahlen ([ˈtsaːlən], "numbers").[4][5]
+
+ℤ is a subset of the set of all rational numbers ℚ, in turn a subset of the real numbers ℝ.  Like the natural numbers, ℤ is countably infinite.
+
+The integers form the smallest group and the smallest ring containing the natural numbers. In algebraic number theory, the integers are sometimes qualified as rational integers to distinguish them from the more general algebraic integers. In fact, the (rational) integers are the algebraic integers that are also rational numbers.
+
+The symbol ℤ can be annotated to denote various sets, with varying usage amongst different authors: ℤ+, ℤ+ or ℤ> for the positive integers, ℤ0+ or ℤ≥ for non-negative integers, ℤ≠ for non-zero integers. Some authors use ℤ* for non-zero integers, others use it for non-negative integers, or for {–1, 1}. Additionally, ℤp is used to denote either the set of integers modulo p, i.e., a set of congruence classes of integers, or the set of p-adic integers.[6][7][8]
+
+Ring homomorphisms
+
+Algebraic structures
+
+Related structures
+
+Algebraic number theory
+
+p-adic number theory and decimals
+
+Algebraic geometry
+
+Noncommutative algebraic geometry
+
+Free algebra
+
+Clifford algebra
+
+Like the natural numbers, ℤ is closed under the operations of addition and multiplication, that is, the sum and product of any two integers is an integer. However, with the inclusion of the negative natural numbers, and, importantly, 0, ℤ (unlike the natural numbers) is also closed under subtraction. The integers form a unital ring which is the most basic one, in the following sense: for any unital ring, there is a unique ring homomorphism from the integers into this ring. This universal property, namely to be an initial object in the category of rings, characterizes the ring ℤ.
+
+ℤ is not closed under division, since the quotient of two integers (e.g., 1 divided by 2), need not be an integer. Although the natural numbers are closed under exponentiation, the integers are not (since the result can be a fraction when the exponent is negative).
+
+The following table lists some of the basic properties of addition and multiplication for any integers a, b and c.
+
+In the language of abstract algebra, the first five properties listed above for addition say that ℤ under addition is an abelian group. It is also a cyclic group, since every non-zero integer can be written as a finite sum 1 + 1 + … + 1 or (−1) + (−1) + … + (−1). In fact, ℤ under addition is the only infinite cyclic group, in the sense that any infinite cyclic group is isomorphic to ℤ.
+
+The first four properties listed above for multiplication say that ℤ under multiplication is a commutative monoid. However, not every integer has a multiplicative inverse; e.g., there is no integer x such that 2x = 1. This means that ℤ under multiplication is not a group.
+
+All the rules from the above property table, except for the last, taken together say that ℤ together with addition and multiplication is a commutative ring with unity. It is the prototype of all objects of such algebraic structure. Only those equalities of expressions are true in ℤ for all values of variables, which are true in any unital commutative ring. Certain non-zero integers map to zero in certain rings.
+
+The lack of zero divisors in the integers (last property in the table) means that the commutative ring ℤ is an integral domain.
+
+The lack of multiplicative inverses, which is equivalent to the fact that ℤ is not closed under division, means that ℤ is not a field. The smallest field containing the integers as a subring is the field of rational numbers. The process of constructing the rationals from the integers can be mimicked to form the field of fractions of any integral domain. And back, starting from an algebraic number field (an extension of rational numbers), its ring of integers can be extracted, which includes ℤ as its subring.
+
+Although ordinary division is not defined on ℤ, the division "with remainder" is defined on them. It is called Euclidean division and possesses the following important property: that is, given two integers a and b with b ≠ 0, there exist unique integers q and r such that a = q × b + r and 0 ≤ r <  | b |, where | b | denotes the absolute value of b. The integer q is called the quotient and r is called the remainder of the division of a by b. The Euclidean algorithm for computing greatest common divisors works by a sequence of Euclidean divisions.
+
+Again, in the language of abstract algebra, the above says that ℤ is a Euclidean domain. This implies that ℤ is a principal ideal domain and any positive integer can be written as the products of primes in an essentially unique way.[9] This is the fundamental theorem of arithmetic.
+
+ℤ is a totally ordered set without upper or lower bound. The ordering of ℤ is given by:
+:... −3 < −2 < −1 < 0 < 1 < 2 < 3 < ...
+An integer is positive if it is greater than zero and negative if it is less than zero. Zero is defined as neither negative nor positive.
+
+The ordering of integers is compatible with the algebraic operations in the following way:
+
+It follows that ℤ together with the above ordering is an ordered ring.
+
+The integers are the only nontrivial totally ordered abelian group whose positive elements are well-ordered.[10] This is equivalent to the statement that any Noetherian valuation ring is either a field or a discrete valuation ring.
+
+In elementary school teaching, integers are often intuitively defined as the (positive) natural numbers, zero, and the negations of the natural numbers. However, this style of definition leads to many different cases (each arithmetic operation needs to be defined on each combination of types of integer) and makes it tedious to prove that these operations obey the laws of arithmetic.[11] Therefore, in modern set-theoretic mathematics a more abstract construction,[12] which allows one to define the arithmetical operations without any case distinction, is often used instead.[13] The integers can thus be formally constructed as the equivalence classes of ordered pairs of natural numbers (a,b).[14]
+
+The intuition is that (a,b) stands for the result of subtracting b from a.[14] To confirm our expectation that 1 − 2 and 4 − 5 denote the same number, we define an equivalence relation ~ on these pairs with the following rule:
+
+precisely when
+
+Addition and multiplication of integers can be defined in terms of the equivalent operations on the natural numbers;[14] denoting by [(a,b)] the equivalence class having (a,b) as a member, one has:
+
+The negation (or additive inverse) of an integer is obtained by reversing the order of the pair:
+
+Hence subtraction can be defined as the addition of the additive inverse:
+
+The standard ordering on the integers is given by:
+
+It is easily verified that these definitions are independent of the choice of representatives of the equivalence classes.
+
+Every equivalence class has a unique member that is of the form (n,0) or (0,n) (or both at once). The natural number n is identified with the class [(n,0)] (in other words the natural numbers are embedded into the integers by map sending n to [(n,0)]), and the class [(0,n)] is denoted −n (this covers all remaining classes, and gives the class [(0,0)] a second time since −0 = 0.
+
+Thus, [(a,b)] is denoted by
+
+If the natural numbers are identified with the corresponding integers (using the embedding mentioned above), this convention creates no ambiguity.
+
+This notation recovers the familiar representation of the integers as  {…, −2, −1, 0, 1, 2, …}.
+
+Some examples are:
+
+In theoretical computer science, other approaches for the construction of integers are used by automated theorem provers and term rewrite engines.
+Integers are represented as algebraic terms built using a few basic operations  (such as zero, succ, pred, etc.) and, possibly, using natural numbers, which are assumed to be already constructed (e.g., using the Peano approach).
+
+There exist at least ten such constructions of signed integers.[15] These constructions differ in several ways: the number of basic operations used for the construction, the number (usually, between 0 and 2) and the types of arguments accepted by these operations; the presence or absence of natural numbers as arguments of some of these operations, and the fact that these operations are free constructors or not, i.e., that the same integer can be represented using only one or many algebraic terms.
+
+The technique for the construction of integers presented above in this section corresponds to the particular case where there is a single basic operation pair
+
+
+
+(
+x
+,
+y
+)
+
+
+{\displaystyle (x,y)}
+
+ that takes as arguments two natural numbers 
+
+
+
+x
+
+
+{\displaystyle x}
+
+ and 
+
+
+
+y
+
+
+{\displaystyle y}
+
+, and returns an integer (equal to 
+
+
+
+x
+−
+y
+
+
+{\displaystyle x-y}
+
+). This operation is not free since the integer 0 can be written pair(0,0), or pair(1,1), or pair(2,2), etc. This technique of construction is used by the proof assistant Isabelle; however, many other tools use alternative construction techniques, notable those based upon free constructors, which are simpler and can be implemented more efficiently in computers.
+
+An integer is often a primitive data type in computer languages. However, integer data types can only represent a subset of all integers, since practical computers are of finite capacity. Also, in the common two's complement representation, the inherent definition of sign distinguishes between "negative" and "non-negative" rather than "negative, positive, and 0". (It is, however, certainly possible for a computer to determine whether an integer value is truly positive.) Fixed length integer approximation data types (or subsets) are denoted int or Integer in several programming languages (such as Algol68, C, Java, Delphi, etc.).
+
+Variable-length representations of integers, such as bignums, can store any integer that fits in the computer's memory. Other integer data types are implemented with a fixed size, usually a number of bits which is a power of 2 (4, 8, 16, etc.) or a memorable number of decimal digits (e.g., 9 or 10).
+
+The cardinality of the set of integers is equal to ℵ0 (aleph-null). This is readily demonstrated by the construction of a bijection, that is, a function that is injective and surjective from ℤ to ℕ.
+If ℕ₀ ≡ {0, 1, 2, ...} then consider the function:
+
+{… (−4,8) (−3,6) (−2,4) (−1,2) (0,0) (1,1) (2,3) (3,5) ...}
+
+If ℕ ≡ {1, 2, 3, ...} then consider the function:
+
+{... (−4,8) (−3,6) (−2,4) (−1,2) (0,1) (1,3) (2,5) (3,7) ...}
+
+If the domain is restricted to ℤ then each and every member of ℤ has one and only one corresponding member of ℕ and by the definition of cardinal equality the two sets have equal cardinality.
+
+This article incorporates material from Integer on PlanetMath, which is licensed under the Creative Commons Attribution/Share-Alike License.

en/1771.html.txt

@@ -0,0 +1 @@
+Set, The Set, or SET may refer to:

en/1772.html.txt

@@ -0,0 +1,162 @@
+
+
+An integer (from the Latin integer meaning "whole")[a] is colloquially defined as a number that can be written without a fractional component. For example, 21, 4, 0, and −2048 are integers, while 9.75, 5+1/2, and √2 are not.
+
+The set of integers consists of zero (0), the positive natural numbers (1, 2, 3, ...), also called whole numbers or counting numbers,[2][3] and their additive inverses (the negative integers, i.e., −1, −2, −3, ...). The set of integers is often denoted by a boldface letter ‘Z’ ("Z") or blackboard bold 
+
+
+
+
+Z
+
+
+
+{\displaystyle \mathbb {Z} }
+
+ (Unicode U+2124 ℤ) standing for the German word Zahlen ([ˈtsaːlən], "numbers").[4][5]
+
+ℤ is a subset of the set of all rational numbers ℚ, in turn a subset of the real numbers ℝ.  Like the natural numbers, ℤ is countably infinite.
+
+The integers form the smallest group and the smallest ring containing the natural numbers. In algebraic number theory, the integers are sometimes qualified as rational integers to distinguish them from the more general algebraic integers. In fact, the (rational) integers are the algebraic integers that are also rational numbers.
+
+The symbol ℤ can be annotated to denote various sets, with varying usage amongst different authors: ℤ+, ℤ+ or ℤ> for the positive integers, ℤ0+ or ℤ≥ for non-negative integers, ℤ≠ for non-zero integers. Some authors use ℤ* for non-zero integers, others use it for non-negative integers, or for {–1, 1}. Additionally, ℤp is used to denote either the set of integers modulo p, i.e., a set of congruence classes of integers, or the set of p-adic integers.[6][7][8]
+
+Ring homomorphisms
+
+Algebraic structures
+
+Related structures
+
+Algebraic number theory
+
+p-adic number theory and decimals
+
+Algebraic geometry
+
+Noncommutative algebraic geometry
+
+Free algebra
+
+Clifford algebra
+
+Like the natural numbers, ℤ is closed under the operations of addition and multiplication, that is, the sum and product of any two integers is an integer. However, with the inclusion of the negative natural numbers, and, importantly, 0, ℤ (unlike the natural numbers) is also closed under subtraction. The integers form a unital ring which is the most basic one, in the following sense: for any unital ring, there is a unique ring homomorphism from the integers into this ring. This universal property, namely to be an initial object in the category of rings, characterizes the ring ℤ.
+
+ℤ is not closed under division, since the quotient of two integers (e.g., 1 divided by 2), need not be an integer. Although the natural numbers are closed under exponentiation, the integers are not (since the result can be a fraction when the exponent is negative).
+
+The following table lists some of the basic properties of addition and multiplication for any integers a, b and c.
+
+In the language of abstract algebra, the first five properties listed above for addition say that ℤ under addition is an abelian group. It is also a cyclic group, since every non-zero integer can be written as a finite sum 1 + 1 + … + 1 or (−1) + (−1) + … + (−1). In fact, ℤ under addition is the only infinite cyclic group, in the sense that any infinite cyclic group is isomorphic to ℤ.
+
+The first four properties listed above for multiplication say that ℤ under multiplication is a commutative monoid. However, not every integer has a multiplicative inverse; e.g., there is no integer x such that 2x = 1. This means that ℤ under multiplication is not a group.
+
+All the rules from the above property table, except for the last, taken together say that ℤ together with addition and multiplication is a commutative ring with unity. It is the prototype of all objects of such algebraic structure. Only those equalities of expressions are true in ℤ for all values of variables, which are true in any unital commutative ring. Certain non-zero integers map to zero in certain rings.
+
+The lack of zero divisors in the integers (last property in the table) means that the commutative ring ℤ is an integral domain.
+
+The lack of multiplicative inverses, which is equivalent to the fact that ℤ is not closed under division, means that ℤ is not a field. The smallest field containing the integers as a subring is the field of rational numbers. The process of constructing the rationals from the integers can be mimicked to form the field of fractions of any integral domain. And back, starting from an algebraic number field (an extension of rational numbers), its ring of integers can be extracted, which includes ℤ as its subring.
+
+Although ordinary division is not defined on ℤ, the division "with remainder" is defined on them. It is called Euclidean division and possesses the following important property: that is, given two integers a and b with b ≠ 0, there exist unique integers q and r such that a = q × b + r and 0 ≤ r <  | b |, where | b | denotes the absolute value of b. The integer q is called the quotient and r is called the remainder of the division of a by b. The Euclidean algorithm for computing greatest common divisors works by a sequence of Euclidean divisions.
+
+Again, in the language of abstract algebra, the above says that ℤ is a Euclidean domain. This implies that ℤ is a principal ideal domain and any positive integer can be written as the products of primes in an essentially unique way.[9] This is the fundamental theorem of arithmetic.
+
+ℤ is a totally ordered set without upper or lower bound. The ordering of ℤ is given by:
+:... −3 < −2 < −1 < 0 < 1 < 2 < 3 < ...
+An integer is positive if it is greater than zero and negative if it is less than zero. Zero is defined as neither negative nor positive.
+
+The ordering of integers is compatible with the algebraic operations in the following way:
+
+It follows that ℤ together with the above ordering is an ordered ring.
+
+The integers are the only nontrivial totally ordered abelian group whose positive elements are well-ordered.[10] This is equivalent to the statement that any Noetherian valuation ring is either a field or a discrete valuation ring.
+
+In elementary school teaching, integers are often intuitively defined as the (positive) natural numbers, zero, and the negations of the natural numbers. However, this style of definition leads to many different cases (each arithmetic operation needs to be defined on each combination of types of integer) and makes it tedious to prove that these operations obey the laws of arithmetic.[11] Therefore, in modern set-theoretic mathematics a more abstract construction,[12] which allows one to define the arithmetical operations without any case distinction, is often used instead.[13] The integers can thus be formally constructed as the equivalence classes of ordered pairs of natural numbers (a,b).[14]
+
+The intuition is that (a,b) stands for the result of subtracting b from a.[14] To confirm our expectation that 1 − 2 and 4 − 5 denote the same number, we define an equivalence relation ~ on these pairs with the following rule:
+
+precisely when
+
+Addition and multiplication of integers can be defined in terms of the equivalent operations on the natural numbers;[14] denoting by [(a,b)] the equivalence class having (a,b) as a member, one has:
+
+The negation (or additive inverse) of an integer is obtained by reversing the order of the pair:
+
+Hence subtraction can be defined as the addition of the additive inverse:
+
+The standard ordering on the integers is given by:
+
+It is easily verified that these definitions are independent of the choice of representatives of the equivalence classes.
+
+Every equivalence class has a unique member that is of the form (n,0) or (0,n) (or both at once). The natural number n is identified with the class [(n,0)] (in other words the natural numbers are embedded into the integers by map sending n to [(n,0)]), and the class [(0,n)] is denoted −n (this covers all remaining classes, and gives the class [(0,0)] a second time since −0 = 0.
+
+Thus, [(a,b)] is denoted by
+
+If the natural numbers are identified with the corresponding integers (using the embedding mentioned above), this convention creates no ambiguity.
+
+This notation recovers the familiar representation of the integers as  {…, −2, −1, 0, 1, 2, …}.
+
+Some examples are:
+
+In theoretical computer science, other approaches for the construction of integers are used by automated theorem provers and term rewrite engines.
+Integers are represented as algebraic terms built using a few basic operations  (such as zero, succ, pred, etc.) and, possibly, using natural numbers, which are assumed to be already constructed (e.g., using the Peano approach).
+
+There exist at least ten such constructions of signed integers.[15] These constructions differ in several ways: the number of basic operations used for the construction, the number (usually, between 0 and 2) and the types of arguments accepted by these operations; the presence or absence of natural numbers as arguments of some of these operations, and the fact that these operations are free constructors or not, i.e., that the same integer can be represented using only one or many algebraic terms.
+
+The technique for the construction of integers presented above in this section corresponds to the particular case where there is a single basic operation pair
+
+
+
+(
+x
+,
+y
+)
+
+
+{\displaystyle (x,y)}
+
+ that takes as arguments two natural numbers 
+
+
+
+x
+
+
+{\displaystyle x}
+
+ and 
+
+
+
+y
+
+
+{\displaystyle y}
+
+, and returns an integer (equal to 
+
+
+
+x
+−
+y
+
+
+{\displaystyle x-y}
+
+). This operation is not free since the integer 0 can be written pair(0,0), or pair(1,1), or pair(2,2), etc. This technique of construction is used by the proof assistant Isabelle; however, many other tools use alternative construction techniques, notable those based upon free constructors, which are simpler and can be implemented more efficiently in computers.
+
+An integer is often a primitive data type in computer languages. However, integer data types can only represent a subset of all integers, since practical computers are of finite capacity. Also, in the common two's complement representation, the inherent definition of sign distinguishes between "negative" and "non-negative" rather than "negative, positive, and 0". (It is, however, certainly possible for a computer to determine whether an integer value is truly positive.) Fixed length integer approximation data types (or subsets) are denoted int or Integer in several programming languages (such as Algol68, C, Java, Delphi, etc.).
+
+Variable-length representations of integers, such as bignums, can store any integer that fits in the computer's memory. Other integer data types are implemented with a fixed size, usually a number of bits which is a power of 2 (4, 8, 16, etc.) or a memorable number of decimal digits (e.g., 9 or 10).
+
+The cardinality of the set of integers is equal to ℵ0 (aleph-null). This is readily demonstrated by the construction of a bijection, that is, a function that is injective and surjective from ℤ to ℕ.
+If ℕ₀ ≡ {0, 1, 2, ...} then consider the function:
+
+{… (−4,8) (−3,6) (−2,4) (−1,2) (0,0) (1,1) (2,3) (3,5) ...}
+
+If ℕ ≡ {1, 2, 3, ...} then consider the function:
+
+{... (−4,8) (−3,6) (−2,4) (−1,2) (0,1) (1,3) (2,5) (3,7) ...}
+
+If the domain is restricted to ℤ then each and every member of ℤ has one and only one corresponding member of ℕ and by the definition of cardinal equality the two sets have equal cardinality.
+
+This article incorporates material from Integer on PlanetMath, which is licensed under the Creative Commons Attribution/Share-Alike License.

en/1773.html.txt

@@ -0,0 +1,162 @@
+
+
+An integer (from the Latin integer meaning "whole")[a] is colloquially defined as a number that can be written without a fractional component. For example, 21, 4, 0, and −2048 are integers, while 9.75, 5+1/2, and √2 are not.
+
+The set of integers consists of zero (0), the positive natural numbers (1, 2, 3, ...), also called whole numbers or counting numbers,[2][3] and their additive inverses (the negative integers, i.e., −1, −2, −3, ...). The set of integers is often denoted by a boldface letter ‘Z’ ("Z") or blackboard bold 
+
+
+
+
+Z
+
+
+
+{\displaystyle \mathbb {Z} }
+
+ (Unicode U+2124 ℤ) standing for the German word Zahlen ([ˈtsaːlən], "numbers").[4][5]
+
+ℤ is a subset of the set of all rational numbers ℚ, in turn a subset of the real numbers ℝ.  Like the natural numbers, ℤ is countably infinite.
+
+The integers form the smallest group and the smallest ring containing the natural numbers. In algebraic number theory, the integers are sometimes qualified as rational integers to distinguish them from the more general algebraic integers. In fact, the (rational) integers are the algebraic integers that are also rational numbers.
+
+The symbol ℤ can be annotated to denote various sets, with varying usage amongst different authors: ℤ+, ℤ+ or ℤ> for the positive integers, ℤ0+ or ℤ≥ for non-negative integers, ℤ≠ for non-zero integers. Some authors use ℤ* for non-zero integers, others use it for non-negative integers, or for {–1, 1}. Additionally, ℤp is used to denote either the set of integers modulo p, i.e., a set of congruence classes of integers, or the set of p-adic integers.[6][7][8]
+
+Ring homomorphisms
+
+Algebraic structures
+
+Related structures
+
+Algebraic number theory
+
+p-adic number theory and decimals
+
+Algebraic geometry
+
+Noncommutative algebraic geometry
+
+Free algebra
+
+Clifford algebra
+
+Like the natural numbers, ℤ is closed under the operations of addition and multiplication, that is, the sum and product of any two integers is an integer. However, with the inclusion of the negative natural numbers, and, importantly, 0, ℤ (unlike the natural numbers) is also closed under subtraction. The integers form a unital ring which is the most basic one, in the following sense: for any unital ring, there is a unique ring homomorphism from the integers into this ring. This universal property, namely to be an initial object in the category of rings, characterizes the ring ℤ.
+
+ℤ is not closed under division, since the quotient of two integers (e.g., 1 divided by 2), need not be an integer. Although the natural numbers are closed under exponentiation, the integers are not (since the result can be a fraction when the exponent is negative).
+
+The following table lists some of the basic properties of addition and multiplication for any integers a, b and c.
+
+In the language of abstract algebra, the first five properties listed above for addition say that ℤ under addition is an abelian group. It is also a cyclic group, since every non-zero integer can be written as a finite sum 1 + 1 + … + 1 or (−1) + (−1) + … + (−1). In fact, ℤ under addition is the only infinite cyclic group, in the sense that any infinite cyclic group is isomorphic to ℤ.
+
+The first four properties listed above for multiplication say that ℤ under multiplication is a commutative monoid. However, not every integer has a multiplicative inverse; e.g., there is no integer x such that 2x = 1. This means that ℤ under multiplication is not a group.
+
+All the rules from the above property table, except for the last, taken together say that ℤ together with addition and multiplication is a commutative ring with unity. It is the prototype of all objects of such algebraic structure. Only those equalities of expressions are true in ℤ for all values of variables, which are true in any unital commutative ring. Certain non-zero integers map to zero in certain rings.
+
+The lack of zero divisors in the integers (last property in the table) means that the commutative ring ℤ is an integral domain.
+
+The lack of multiplicative inverses, which is equivalent to the fact that ℤ is not closed under division, means that ℤ is not a field. The smallest field containing the integers as a subring is the field of rational numbers. The process of constructing the rationals from the integers can be mimicked to form the field of fractions of any integral domain. And back, starting from an algebraic number field (an extension of rational numbers), its ring of integers can be extracted, which includes ℤ as its subring.
+
+Although ordinary division is not defined on ℤ, the division "with remainder" is defined on them. It is called Euclidean division and possesses the following important property: that is, given two integers a and b with b ≠ 0, there exist unique integers q and r such that a = q × b + r and 0 ≤ r <  | b |, where | b | denotes the absolute value of b. The integer q is called the quotient and r is called the remainder of the division of a by b. The Euclidean algorithm for computing greatest common divisors works by a sequence of Euclidean divisions.
+
+Again, in the language of abstract algebra, the above says that ℤ is a Euclidean domain. This implies that ℤ is a principal ideal domain and any positive integer can be written as the products of primes in an essentially unique way.[9] This is the fundamental theorem of arithmetic.
+
+ℤ is a totally ordered set without upper or lower bound. The ordering of ℤ is given by:
+:... −3 < −2 < −1 < 0 < 1 < 2 < 3 < ...
+An integer is positive if it is greater than zero and negative if it is less than zero. Zero is defined as neither negative nor positive.
+
+The ordering of integers is compatible with the algebraic operations in the following way:
+
+It follows that ℤ together with the above ordering is an ordered ring.
+
+The integers are the only nontrivial totally ordered abelian group whose positive elements are well-ordered.[10] This is equivalent to the statement that any Noetherian valuation ring is either a field or a discrete valuation ring.
+
+In elementary school teaching, integers are often intuitively defined as the (positive) natural numbers, zero, and the negations of the natural numbers. However, this style of definition leads to many different cases (each arithmetic operation needs to be defined on each combination of types of integer) and makes it tedious to prove that these operations obey the laws of arithmetic.[11] Therefore, in modern set-theoretic mathematics a more abstract construction,[12] which allows one to define the arithmetical operations without any case distinction, is often used instead.[13] The integers can thus be formally constructed as the equivalence classes of ordered pairs of natural numbers (a,b).[14]
+
+The intuition is that (a,b) stands for the result of subtracting b from a.[14] To confirm our expectation that 1 − 2 and 4 − 5 denote the same number, we define an equivalence relation ~ on these pairs with the following rule:
+
+precisely when
+
+Addition and multiplication of integers can be defined in terms of the equivalent operations on the natural numbers;[14] denoting by [(a,b)] the equivalence class having (a,b) as a member, one has:
+
+The negation (or additive inverse) of an integer is obtained by reversing the order of the pair:
+
+Hence subtraction can be defined as the addition of the additive inverse:
+
+The standard ordering on the integers is given by:
+
+It is easily verified that these definitions are independent of the choice of representatives of the equivalence classes.
+
+Every equivalence class has a unique member that is of the form (n,0) or (0,n) (or both at once). The natural number n is identified with the class [(n,0)] (in other words the natural numbers are embedded into the integers by map sending n to [(n,0)]), and the class [(0,n)] is denoted −n (this covers all remaining classes, and gives the class [(0,0)] a second time since −0 = 0.
+
+Thus, [(a,b)] is denoted by
+
+If the natural numbers are identified with the corresponding integers (using the embedding mentioned above), this convention creates no ambiguity.
+
+This notation recovers the familiar representation of the integers as  {…, −2, −1, 0, 1, 2, …}.
+
+Some examples are:
+
+In theoretical computer science, other approaches for the construction of integers are used by automated theorem provers and term rewrite engines.
+Integers are represented as algebraic terms built using a few basic operations  (such as zero, succ, pred, etc.) and, possibly, using natural numbers, which are assumed to be already constructed (e.g., using the Peano approach).
+
+There exist at least ten such constructions of signed integers.[15] These constructions differ in several ways: the number of basic operations used for the construction, the number (usually, between 0 and 2) and the types of arguments accepted by these operations; the presence or absence of natural numbers as arguments of some of these operations, and the fact that these operations are free constructors or not, i.e., that the same integer can be represented using only one or many algebraic terms.
+
+The technique for the construction of integers presented above in this section corresponds to the particular case where there is a single basic operation pair
+
+
+
+(
+x
+,
+y
+)
+
+
+{\displaystyle (x,y)}
+
+ that takes as arguments two natural numbers 
+
+
+
+x
+
+
+{\displaystyle x}
+
+ and 
+
+
+
+y
+
+
+{\displaystyle y}
+
+, and returns an integer (equal to 
+
+
+
+x
+−
+y
+
+
+{\displaystyle x-y}
+
+). This operation is not free since the integer 0 can be written pair(0,0), or pair(1,1), or pair(2,2), etc. This technique of construction is used by the proof assistant Isabelle; however, many other tools use alternative construction techniques, notable those based upon free constructors, which are simpler and can be implemented more efficiently in computers.
+
+An integer is often a primitive data type in computer languages. However, integer data types can only represent a subset of all integers, since practical computers are of finite capacity. Also, in the common two's complement representation, the inherent definition of sign distinguishes between "negative" and "non-negative" rather than "negative, positive, and 0". (It is, however, certainly possible for a computer to determine whether an integer value is truly positive.) Fixed length integer approximation data types (or subsets) are denoted int or Integer in several programming languages (such as Algol68, C, Java, Delphi, etc.).
+
+Variable-length representations of integers, such as bignums, can store any integer that fits in the computer's memory. Other integer data types are implemented with a fixed size, usually a number of bits which is a power of 2 (4, 8, 16, etc.) or a memorable number of decimal digits (e.g., 9 or 10).
+
+The cardinality of the set of integers is equal to ℵ0 (aleph-null). This is readily demonstrated by the construction of a bijection, that is, a function that is injective and surjective from ℤ to ℕ.
+If ℕ₀ ≡ {0, 1, 2, ...} then consider the function:
+
+{… (−4,8) (−3,6) (−2,4) (−1,2) (0,0) (1,1) (2,3) (3,5) ...}
+
+If ℕ ≡ {1, 2, 3, ...} then consider the function:
+
+{... (−4,8) (−3,6) (−2,4) (−1,2) (0,1) (1,3) (2,5) (3,7) ...}
+
+If the domain is restricted to ℤ then each and every member of ℤ has one and only one corresponding member of ℕ and by the definition of cardinal equality the two sets have equal cardinality.
+
+This article incorporates material from Integer on PlanetMath, which is licensed under the Creative Commons Attribution/Share-Alike License.

en/1774.html.txt

@@ -0,0 +1,36 @@
+
+
+A company, abbreviated as co., is a legal entity representing an association of people, whether natural, legal or a mixture of both, with a specific objective. Company members share a common purpose and unite to achieve specific, declared goals. Companies take various forms, such as:
+
+A company can be created as a legal person so that the company itself has limited liability as members perform or fail to discharge their duty according to the publicly declared incorporation, or published policy. When a company closes, it may need to be liquidated to avoid further legal obligations.
+
+Companies may associate and collectively register themselves as new companies; the resulting entities are often known as corporate groups.
+
+One can define a company as an "artificial person", invisible, intangible, created by or under law,[1] with a discrete legal personality, perpetual succession, and a common seal. Except for some senior positions, companies remain unaffected by the death, insanity, or insolvency of an individual member.
+
+The English word company has its origins in the Old French term compagnie (first recorded in 1150), meaning a "society, friendship, intimacy; body of soldiers",[2] which came from the Late Latin word companio ("one who eats bread with you"), first attested in the Lex Salica (English: Salic Law) (c. 500 CE) as a calque of the  Germanic expression gahlaibo (literally, "with bread"), related to Old High German galeipo ("companion") and to Gothic gahlaiba ("messmate").
+
+By 1303, the word referred to trade guilds.[3] Usage of the term company to mean "business association" was first recorded in 1553,[4]
+and the abbreviation "co." dates from 1769.[5][6]
+
+In English law and in legal jurisdictions based upon it, a company is a body corporate or corporation company registered under the Companies Acts or under similar legislation.[7] Common forms include:
+
+In the United Kingdom, a partnership is not legally a company, but may sometimes be referred to (informally) as a company. It may be referred to as a  firm.
+
+In the United States, a company may be a "corporation, partnership, association, joint-stock company,  trust,  fund, or organized group of persons, whether incorporated or not, and (in an official capacity) any receiver, trustee in bankruptcy, or similar official, or liquidating agent, for any of the foregoing".[8][9]  In the US, a company is not necessarily a corporation.[10]
+
+Less common types of companies are:
+
+When "Ltd" is placed after the company's name, it signifies a limited company, and "PLC" (public limited company) indicates that its shares are widely held.[13]
+
+In the legal context, the owners of a company are normally referred to as the "members". In a company limited or unlimited by shares (formed or incorporated with a share capital), this will be the shareholders. In a company limited by guarantee, this will be the guarantors. Some offshore jurisdictions have created special forms of offshore company in a bid to attract business for their jurisdictions. Examples include "segregated portfolio companies" and restricted purpose companies.
+
+There are, however, many, many sub-categories of types of company that can be formed in various jurisdictions in the world.
+
+Companies are also sometimes distinguished for legal and regulatory purposes between public companies and private companies. Public companies are companies whose shares can be publicly traded, often (although not always) on a stock exchange which imposes listing requirements/Listing Rules as to the issued shares, the trading of shares and future issue of shares to help bolster the reputation of the exchange or particular market of an exchange. Private companies do not have publicly traded shares, and often contain restrictions on transfers of shares. In some jurisdictions, private companies have maximum numbers of shareholders.
+
+A parent company is a company that owns enough voting stock in another firm to control management and operations by influencing or electing its board of directors; the second company being deemed as a subsidiary of the parent company. The definition of a parent company differs by jurisdiction, with the definition normally being defined by way of laws dealing with companies in that jurisdiction.
+
+In China companies are often government run or some government supported. Others may be Foreign Companies or export-based corporations. However, many of these companies are government regulated.
+
+Garner, Bryan A., ed. (1891). "company". Black's Law Dictionary. Black's Law, 9th Edition. 1 (9 ed.). St. Paul, Minnesota: West Publishing, Inc (published 2009). p. 318. ISBN 9780314199492. Retrieved April 20, 2019. 2. A corporation, partnership, association, joint-stock company, trust, fund, or organized group of persons, whether incorporated or not, and (in an official capacity) any receiver, trustee in bankruptcy, or similar official, or liquidating agent, for any of the foregoing. Investment Company Act 2(a)(8)(15 USCA 80a-2(a)(8)).

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+
+
+A company, abbreviated as co., is a legal entity representing an association of people, whether natural, legal or a mixture of both, with a specific objective. Company members share a common purpose and unite to achieve specific, declared goals. Companies take various forms, such as:
+
+A company can be created as a legal person so that the company itself has limited liability as members perform or fail to discharge their duty according to the publicly declared incorporation, or published policy. When a company closes, it may need to be liquidated to avoid further legal obligations.
+
+Companies may associate and collectively register themselves as new companies; the resulting entities are often known as corporate groups.
+
+One can define a company as an "artificial person", invisible, intangible, created by or under law,[1] with a discrete legal personality, perpetual succession, and a common seal. Except for some senior positions, companies remain unaffected by the death, insanity, or insolvency of an individual member.
+
+The English word company has its origins in the Old French term compagnie (first recorded in 1150), meaning a "society, friendship, intimacy; body of soldiers",[2] which came from the Late Latin word companio ("one who eats bread with you"), first attested in the Lex Salica (English: Salic Law) (c. 500 CE) as a calque of the  Germanic expression gahlaibo (literally, "with bread"), related to Old High German galeipo ("companion") and to Gothic gahlaiba ("messmate").
+
+By 1303, the word referred to trade guilds.[3] Usage of the term company to mean "business association" was first recorded in 1553,[4]
+and the abbreviation "co." dates from 1769.[5][6]
+
+In English law and in legal jurisdictions based upon it, a company is a body corporate or corporation company registered under the Companies Acts or under similar legislation.[7] Common forms include:
+
+In the United Kingdom, a partnership is not legally a company, but may sometimes be referred to (informally) as a company. It may be referred to as a  firm.
+
+In the United States, a company may be a "corporation, partnership, association, joint-stock company,  trust,  fund, or organized group of persons, whether incorporated or not, and (in an official capacity) any receiver, trustee in bankruptcy, or similar official, or liquidating agent, for any of the foregoing".[8][9]  In the US, a company is not necessarily a corporation.[10]
+
+Less common types of companies are:
+
+When "Ltd" is placed after the company's name, it signifies a limited company, and "PLC" (public limited company) indicates that its shares are widely held.[13]
+
+In the legal context, the owners of a company are normally referred to as the "members". In a company limited or unlimited by shares (formed or incorporated with a share capital), this will be the shareholders. In a company limited by guarantee, this will be the guarantors. Some offshore jurisdictions have created special forms of offshore company in a bid to attract business for their jurisdictions. Examples include "segregated portfolio companies" and restricted purpose companies.
+
+There are, however, many, many sub-categories of types of company that can be formed in various jurisdictions in the world.
+
+Companies are also sometimes distinguished for legal and regulatory purposes between public companies and private companies. Public companies are companies whose shares can be publicly traded, often (although not always) on a stock exchange which imposes listing requirements/Listing Rules as to the issued shares, the trading of shares and future issue of shares to help bolster the reputation of the exchange or particular market of an exchange. Private companies do not have publicly traded shares, and often contain restrictions on transfers of shares. In some jurisdictions, private companies have maximum numbers of shareholders.
+
+A parent company is a company that owns enough voting stock in another firm to control management and operations by influencing or electing its board of directors; the second company being deemed as a subsidiary of the parent company. The definition of a parent company differs by jurisdiction, with the definition normally being defined by way of laws dealing with companies in that jurisdiction.
+
+In China companies are often government run or some government supported. Others may be Foreign Companies or export-based corporations. However, many of these companies are government regulated.
+
+Garner, Bryan A., ed. (1891). "company". Black's Law Dictionary. Black's Law, 9th Edition. 1 (9 ed.). St. Paul, Minnesota: West Publishing, Inc (published 2009). p. 318. ISBN 9780314199492. Retrieved April 20, 2019. 2. A corporation, partnership, association, joint-stock company, trust, fund, or organized group of persons, whether incorporated or not, and (in an official capacity) any receiver, trustee in bankruptcy, or similar official, or liquidating agent, for any of the foregoing. Investment Company Act 2(a)(8)(15 USCA 80a-2(a)(8)).

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+An envelope is a common packaging item, usually made of thin, flat material. It is designed to contain a flat object, such as a letter or card.
+
+Traditional envelopes are made from sheets of paper cut to one of three shapes: a rhombus, a short-arm cross or a kite. These shapes allow for the creation of the envelope structure by folding the sheet sides around a central rectangular area. In this manner, a rectangle-faced enclosure is formed with an arrangement of four flaps on the reverse side.
+
+When the folding sequence is such that the last flap to be closed is on a short side it is referred to in commercial envelope manufacture as a pocket - a format frequently employed in the packaging of small quantities of seeds. Although in principle the flaps can be held in place by securing the topmost flap at a single point (for example with a wax seal), generally they are pasted or gummed together at the overlaps. They are most commonly used for enclosing and sending mail (letters) through a prepaid-postage postal system.
+
+Window envelopes have a hole cut in the front side that allows the paper within to be seen.[1] They are generally arranged so that the receiving address printed on the letter is visible, saving the sender from having to duplicate the address on the envelope itself. The window is normally covered with a transparent or translucent film to protect the letter inside, as was first designed by Americus F. Callahan in 1901 and patented the following year.[2] In some cases, shortages of materials or the need to economize resulted in envelopes that had no film covering the window.[citation needed] One innovative process, invented in Europe about 1905, involved using hot oil to saturate the area of the envelope where the address would appear. The treated area became sufficiently translucent for the address to be readable. As of 2009[update] there is no international standard for window envelopes, but some countries, including Germany and the United Kingdom, have national standards.[3]
+
+An aerogram is related to a lettersheet, both being designed to have writing on the inside to minimize the weight. Any handmade envelope is effectively a lettersheet because prior to the folding stage it offers the opportunity for writing a message on that area of the sheet that after folding becomes the inside of the face of the envelope.
+
+The "envelope" used to launch the Penny Post component of the British postal reforms of 1840 by Sir Rowland Hill and the invention of the postage stamp, was a lozenge-shaped lettersheet known as a Mulready.[4] If desired, a separate letter could be enclosed with postage remaining at one penny provided the combined weight did not exceed half an ounce (14 grams). This was a legacy of the previous system of calculating postage, which partly depended on the number of sheets of paper used.
+
+During the U.S. Civil War those in the Confederate States Army occasionally used envelopes made from wallpaper, due to financial hardship.
+
+A "return envelope" is a pre-addressed, smaller envelope included as the contents of a larger envelope and can be used for courtesy reply mail, metered reply mail, or freepost (business reply mail). Some envelopes are designed to be reused as the return envelope, saving the expense of including a return envelope in the contents of the original envelope. The direct mail industry makes extensive use of return envelopes as a response mechanism.
+
+Up until 1840, all envelopes were handmade, each being individually cut to the appropriate shape out of an individual rectangular sheet. In that year George Wilson in the United Kingdom patented the method of tessellating (tiling) a number of envelope patterns across and down a large sheet, thereby reducing the overall amount of waste produced per envelope when they were cut out. In 1845 Edwin Hill and Warren de la Rue obtained a patent for a steam-driven machine that not only cut out the envelope shapes but creased and folded them as well. (Mechanised gumming had yet to be devised.) The convenience of the sheets ready cut to shape popularized the use of machine-made envelopes, and the economic significance of the factories that had produced handmade envelopes gradually diminished.
+
+As envelopes are made of paper, they are intrinsically amenable to embellishment with additional graphics and text over and above the necessary postal markings. This is a feature that the direct mail industry has long taken advantage of—and more recently the Mail Art movement. Custom printed envelopes has also become an increasingly popular marketing method for small business.
+
+Most of the over 400 billion envelopes of all sizes made worldwide are machine-made.
+
+International standard ISO 269 (since withdrawn) defined several standard envelope sizes, which are designed for use with ISO 216 standard paper sizes:
+
+The German standard DIN 678 defines a similar list of envelope formats.
+
+There are dozens of sizes of envelopes available in the United States.
+
+The designations such as "A2" do not correspond to ISO paper sizes. Sometimes, North American paper jobbers and printers will insert a hyphen to distinguish from ISO sizes, thus: A-2.
+
+The No. 10 envelope is the standard business envelope size in the United States.[5]
+PWG 5101.1[6] also lists the following even inch sizes for envelopes: 6 × 9, 7 × 9, 9 × 11, 9 × 12, 10 × 13, 10 × 14 and 10 × 15.
+
+Envelopes accepted by the U.S. Postal Service for mailing at the price of a letter must be:
+
+Japanese traditional rectangular (kakukei) and long (chōkei) envelopes open on the short side, while Western-style (youkei) envelopes open on the long side.
+
+
+
+The first known envelope was nothing like the paper envelope of today. It can be dated back to around 3500 to 3200 BC in the ancient Middle East. Hollow, clay spheres were molded around financial tokens and used in private transactions. The two people who discovered these first envelopes were Jacques de Morgan, in 1901, and Roland de Mecquenem, in 1907.
+
+Paper envelopes were developed in China, where paper was invented by 2nd century BC.[9] Paper envelopes, known as chih poh, were used to store gifts of money. In the Southern Song dynasty, the Chinese imperial court used paper envelopes to distribute monetary gifts to government officials.[10]
+
+Prior to 1845, hand-made envelopes were all that were available for use, both commercial and domestic. In 1845, Edwin Hill and Warren De La Rue were granted a British patent for the first envelope-making machine.[11]
+
+The "envelopes" produced by the Hill/De La Rue machine were not like those used today. They were flat diamond, lozenge (or rhombus)-shaped sheets or "blanks" that had been precut to shape before being fed to the machine for creasing and made ready for folding to form a rectangular enclosure. The edges of the overlapping flaps treated with a paste or adhesive and the method of securing the envelope or wrapper was a user choice. The symmetrical flap arrangement meant that it could be held together with a single wax seal at the apex of the topmost flap. (That the flaps of an envelope can be held together by applying a seal at a single point is a classic design feature of an envelope.)[citation needed]
+
+Nearly 50 years passed before a commercially successful machine for producing pre-gummed envelopes, like those in use today, appeared.
+
+The origin of the use of the diamond shape for envelopes is debated. However, as an alternative to simply wrapping a sheet of paper around a folded letter or an invitation and sealing the edges, it is a tidy and ostensibly paper-efficient way of producing a rectangular-faced envelope. Where the claim to be paper-efficient fails is a consequence of paper manufacturers normally making paper available in rectangular sheets, because the largest size of envelope that can be realised by cutting out a diamond or any other shape which yields an envelope with symmetrical flaps is smaller than the largest that can be made from that sheet simply by folding.
+
+The folded diamond-shaped sheet (or "blank") was in use at the beginning of the 19th century as a novelty wrapper for invitations and letters among the proportion of the population that had the time to sit and cut them out and were affluent enough not to bother about the waste offcuts.  Their use first became widespread in the UK when the British government took monopoly control of postal services and tasked Rowland Hill with its introduction. The new service was launched in May 1840 with a postage-paid machine-printed illustrated (or pictorial) version of the wrapper and the much-celebrated first adhesive postage stamp, the Penny Black, for the production of which the Jacob Perkins printing process was used to deter counterfeiting and forgery.  The wrappers were printed and sold as a sheet of 12, with cutting the purchaser's task. Known as Mulready stationery, because the illustration was created by the respected artist William Mulready, the envelopes were withdrawn when the illustration was ridiculed and lampooned. Nevertheless, the public apparently saw the convenience of the wrappers being available ready-shaped, and it must have been obvious that with the stamp available totally plain versions of the wrapper could be produced and postage prepaid by purchasing a stamp and affixing it to the wrapper once folded and secured.  In this way although the postage-prepaid printed pictorial version died ignominiously, the diamond-shaped wrapper acquired de facto official status and became readily available to the public notwithstanding the time taken to cut them out and the waste generated. With the issuing of the stamps and the operation and control of the service (which is a communications medium) in government hands the British model spread around the world and the diamond-shaped wrapper went with it.
+
+Hill also installed his brother Edwin as The Controller of Stamps, and it was he with his partner Warren De La Rue who patented the machine for mass-producing the diamond-shaped sheets for conversion to envelopes in 1845. Today, envelope-making machine manufacture is a long- and well-established international industry, and blanks are produced with a short-arm-cross shape and a kite shape as well as diamond shape. (The short-arm-cross style is mostly encountered in "pocket" envelopes i.e. envelopes with the closing flap on a short side. The more common style, with the closing flap on a long side, are sometimes referred to as "standard" or "wallet" style for purposes of differentiation.)
+
+The most famous paper-making machine was the Fourdrinier machine. The process involves taking processed pulp stock and converting it to a continuous web which is gathered as a reel.  Subsequently, the reel is guillotined edge to edge to create a large number of properly rectangular sheets because ever since the invention of Gutenberg's press paper has been closely associated with printing.
+
+To this day, all other mechanical printing and duplicating equipments devised in the meantime, including the typewriter (which was used up to the 1990s for addressing envelopes), have been primarily designed to process rectangular sheets. Hence the large sheets are in turn are guillotined down to the sizes of rectangular sheet commonly used in the commercial printing industry, and nowadays to the sizes commonly used as feed-stock in office-grade computer printers, copiers and duplicators (mainly ISO, A4 and US Letter).
+
+Using any mechanical printing equipment to print on envelopes, which although rectangular, are in fact folded sheets with differing thicknesses across their surfaces, calls for skill and attention on the part of the operator. In commercial printing the task of printing on machine-made envelopes is referred to as "overprinting" and is usually confined to the front of the envelope. If printing is required on all four flaps as well as the front, the process is referred to as "printing on the flat". Eye-catching illustrated envelopes or pictorial envelopes, the origins of which as an artistic genre can be attributed to the Mulready stationery – and which was printed in this way - are used extensively for direct mail. In this respect,  direct mail envelopes have a shared history with propaganda envelopes (or "covers") as they are called by philatelists.
+
+At the end of the 20th century, in 1998, the digital printing revolution delivered another benefit for small businesses when the U.S. Postal Service became the first postal authority to approve the introduction of a system of applying to an envelope in the printer bin of a PC sheet printer a digital frank or stamp delivered via the Internet. With this innovative alternative to an adhesive-backed postage stamp as the basis for an  Electronic Stamp Distribution (ESD) service, a business envelope could be produced in-house, addressed and customized with advertising information on the face, and ready to be mailed.
+
+The fortunes of the commercial envelope manufacturing industry and the postal service go hand in hand, and both link to the printing industry and the mechanized envelope processing industry producing equipment such as franking and addressing machines. They are all four symbiotic: technological developments affecting one obviously ricochet through the others: addressing machines print addresses, postage stamps are a print product, franking machines imprint a frank on an envelope. If fewer envelopes are required; fewer stamps are required; fewer franking machines are required and fewer addressing machines are required.[citation needed] For example, the advent and adoption of information-based indicia (IBI) (commonly referred to as digitally-encoded electronic stamps or digital indicia) by the US Postal Service in 1998 caused widespread consternation in the franking machine industry, as their equipments were effectively rendered obsolescent and resulted in a flurry of lawsuits involving Pitney Bowes among others. The advent of e-mail in the late 1990s appeared to offer a substantial threat to the postal service. By 2008 letter-post service operators were reporting significantly smaller volumes of letter-post, specifically stamped envelopes, which they attributed mainly to replacement by e-mail. Although a corresponding reduction in the volume of envelopes required would have been expected, no such decrease was reported as widely as the reduction in letter-post volumes.
+
+Although, with regards to e-mail developments, there is a substantial threat of "technology replacing tradition". This is offset by the equal reasoning that the Universal Postal Union is an international specialised agency of the United Nations, and a source of revenue for the government. Consequently, any deterioration of domestic and international postal services attended by loss of revenue is a matter of governmental concern.
+
+A windowed envelope is an envelope with a plastic or glassine window in it. The plastic in these envelopes creates problems in paper recycling.
+
+Security envelopes have special tamper-resistant and tamper-evident features.  They are used for high value products and documents as well as for evidence for legal proceedings.
+
+Some security envelopes have a patterned tint printed on the inside, which makes it difficult to read the contents. Various patterns exist.[12]
+
+Some envelopes are available for full-size documents or for other items. Some carriers have large mailing envelopes for their express services. Other similar envelopes are available at stationery supply locations.
+
+These mailers usually have an opening on an end with a flap that can be attached by gummed adhesive, integral pressure-sensitive adhesive, adhesive tape, or security tape.
+Construction is usually:
+
+Shipping envelopes can have padding to provide stiffness and some degree of cushioning.  The padding can be ground newsprint, plastic foam sheets, or bubble packing.
+
+Various U.S. Federal Government offices use Standard Form (SF) 65 Government Messenger Envelopes for inter-office mail delivery.  These envelopes are typically light brown in color and un-sealed with string-tied closure method and an array of holes throughout both sides such that it is somewhat visible what the envelope contains.  Other colloquial names for this envelope include "Holey Joe" and "Shotgun" envelope due to the holey nature of the envelope.  Address method is unique in that these envelopes are re-usable and the previous address is crossed out thoroughly and the new addressee (name, building, room, and mailstop) is written in the next available box.  Although still in use, SF-65 is no longer listed on the United States Office of Personnel Management website list of standard forms, which may indicate new envelopes are no longer being printed.  [13]

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+Environment most often refers to:

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+
+
+Enzymes /ˈɛnzaɪmz/ are proteins that act as biological catalysts (biocatalysts). Catalysts accelerate chemical reactions. The molecules upon which enzymes may act are called substrates, and the enzyme converts the substrates into different molecules known as products. Almost all metabolic processes in the cell need enzyme catalysis in order to occur at rates fast enough to sustain life.[1]:8.1 Metabolic pathways depend upon enzymes to catalyze individual steps. The study of enzymes is called enzymology and a new field of pseudoenzyme analysis has recently grown up, recognising that during evolution, some enzymes have lost the ability to carry out biological catalysis, which is often reflected in their amino acid sequences and unusual 'pseudocatalytic' properties.[2][3]
+
+Enzymes are known to catalyze more than 5,000 biochemical reaction types.[4] Other biocatalysts are catalytic RNA molecules, called ribozymes. Enzymes' specificity comes from their unique three-dimensional structures.
+
+Like all catalysts, enzymes increase the reaction rate by lowering its activation energy. Some enzymes can make their conversion of substrate to product occur many millions of times faster. An extreme example is orotidine 5'-phosphate decarboxylase, which allows a reaction that would otherwise take millions of years to occur in milliseconds.[5][6] Chemically, enzymes are like any catalyst and are not consumed in chemical reactions, nor do they alter the equilibrium of a reaction. Enzymes differ from most other catalysts by being much more specific. Enzyme activity can be affected by other molecules: inhibitors are molecules that decrease enzyme activity, and activators are molecules that increase activity. Many therapeutic drugs and poisons are enzyme inhibitors. An enzyme's activity decreases markedly outside its optimal temperature and pH, and many enzymes are (permanently) denatured when exposed to excessive heat, losing their structure and catalytic properties.
+
+Some enzymes are used commercially, for example, in the synthesis of antibiotics. Some household products use enzymes to speed up chemical reactions: enzymes in biological washing powders break down protein, starch or fat stains on clothes, and enzymes in meat tenderizer break down proteins into smaller molecules, making the meat easier to chew.
+
+By the late 17th and early 18th centuries, the digestion of meat by stomach secretions[7] and the conversion of starch to sugars by plant extracts and saliva were known but the mechanisms by which these occurred had not been identified.[8]
+
+French chemist Anselme Payen was the first to discover an enzyme, diastase, in 1833.[9] A few decades later, when studying the fermentation of sugar to alcohol by yeast, Louis Pasteur concluded that this fermentation was caused by a vital force contained within the yeast cells called "ferments", which were thought to function only within living organisms. He wrote that "alcoholic fermentation is an act correlated with the life and organization of the yeast cells, not with the death or putrefaction of the cells."[10]
+
+In 1877, German physiologist Wilhelm Kühne (1837–1900) first used the term enzyme, which comes from Greek ἔνζυμον, "leavened" or "in yeast", to describe this process.[11] The word enzyme was used later to refer to nonliving substances such as pepsin, and the word ferment was used to refer to chemical activity produced by living organisms.[12]
+
+Eduard Buchner submitted his first paper on the study of yeast extracts in 1897. In a series of experiments at the University of Berlin, he found that sugar was fermented by yeast extracts even when there were no living yeast cells in the mixture.[13] He named the enzyme that brought about the fermentation of sucrose "zymase".[14] In 1907, he received the Nobel Prize in Chemistry for "his discovery of cell-free fermentation". Following Buchner's example, enzymes are usually named according to the reaction they carry out: the suffix -ase is combined with the name of the substrate (e.g., lactase is the enzyme that cleaves lactose) or to the type of reaction (e.g., DNA polymerase forms DNA polymers).[15]
+
+The biochemical identity of enzymes was still unknown in the early 1900s. Many scientists observed that enzymatic activity was associated with proteins, but others (such as Nobel laureate Richard Willstätter) argued that proteins were merely carriers for the true enzymes and that proteins per se were incapable of catalysis.[16] In 1926, James B. Sumner showed that the enzyme urease was a pure protein and crystallized it; he did likewise for the enzyme catalase in 1937. The conclusion that pure proteins can be enzymes was definitively demonstrated by John Howard Northrop and Wendell Meredith Stanley, who worked on the digestive enzymes pepsin (1930), trypsin and chymotrypsin. These three scientists were awarded the 1946 Nobel Prize in Chemistry.[17]
+
+The discovery that enzymes could be crystallized eventually allowed their structures to be solved by x-ray crystallography. This was first done for lysozyme, an enzyme found in tears, saliva and egg whites that digests the coating of some bacteria; the structure was solved by a group led by David Chilton Phillips and published in 1965.[18] This high-resolution structure of lysozyme marked the beginning of the field of structural biology and the effort to understand how enzymes work at an atomic level of detail.[19]
+
+An enzyme's name is often derived from its substrate or the chemical reaction it catalyzes, with the word ending in -ase.[1]:8.1.3 Examples are lactase, alcohol dehydrogenase and DNA polymerase. Different enzymes that catalyze the same chemical reaction are called isozymes.[1]:10.3
+
+The International Union of Biochemistry and Molecular Biology have developed a nomenclature for enzymes, the EC numbers; each enzyme is described by a sequence of four numbers preceded by "EC", which stands for "Enzyme Commission". The first number broadly classifies the enzyme based on its mechanism.[20]
+
+The top-level classification is:
+
+These sections are subdivided by other features such as the substrate, products, and chemical mechanism. An enzyme is fully specified by four numerical designations. For example, hexokinase (EC 2.7.1.1) is a transferase (EC 2) that adds a phosphate group (EC 2.7) to a hexose sugar, a molecule containing an alcohol group (EC 2.7.1).[21]
+
+Enzymes are generally globular proteins, acting alone or in larger complexes. The sequence of the amino acids specifies the structure which in turn determines the catalytic activity of the enzyme.[22] Although structure determines function, a novel enzymatic activity cannot yet be predicted from structure alone.[23] Enzyme structures unfold (denature) when heated or exposed to chemical denaturants and this disruption to the structure typically causes a loss of activity.[24] Enzyme denaturation is normally linked to temperatures above a species' normal level; as a result, enzymes from bacteria living in volcanic environments such as hot springs are prized by industrial users for their ability to function at high temperatures, allowing enzyme-catalysed reactions to be operated at a very high rate.
+
+Enzymes are usually much larger than their substrates. Sizes range from just 62 amino acid residues, for the monomer of 4-oxalocrotonate tautomerase,[25] to over 2,500 residues in the animal fatty acid synthase.[26] Only a small portion of their structure (around 2–4 amino acids) is directly involved in catalysis: the catalytic site.[27] This catalytic site is located next to one or more binding sites where residues orient the substrates. The catalytic site and binding site together compose the enzyme's active site. The remaining majority of the enzyme structure serves to maintain the precise orientation and dynamics of the active site.[28]
+
+In some enzymes, no amino acids are directly involved in catalysis; instead, the enzyme contains sites to bind and orient catalytic cofactors.[28] Enzyme structures may also contain allosteric sites where the binding of a small molecule causes a conformational change that increases or decreases activity.[29]
+
+A small number of RNA-based biological catalysts called ribozymes exist, which again can act alone or in complex with proteins. The most common of these is the ribosome which is a complex of protein and catalytic RNA components.[1]:2.2
+
+Enzymes must bind their substrates before they can catalyse any chemical reaction. Enzymes are usually very specific as to what substrates they bind and then the chemical reaction catalysed. Specificity is achieved by binding pockets with complementary shape, charge and hydrophilic/hydrophobic characteristics to the substrates. Enzymes can therefore distinguish between very similar substrate molecules to be chemoselective, regioselective and stereospecific.[30]
+
+Some of the enzymes showing the highest specificity and accuracy are involved in the copying and expression of the genome. Some of these enzymes have "proof-reading" mechanisms. Here, an enzyme such as DNA polymerase catalyzes a reaction in a first step and then checks that the product is correct in a second step.[31] This two-step process results in average error rates of less than 1 error in 100 million reactions in high-fidelity mammalian polymerases.[1]:5.3.1 Similar proofreading mechanisms are also found in RNA polymerase,[32] aminoacyl tRNA synthetases[33] and ribosomes.[34]
+
+Conversely, some enzymes display enzyme promiscuity, having broad specificity and acting on a range of different physiologically relevant substrates. Many enzymes possess small side activities which arose fortuitously (i.e. neutrally), which may be the starting point for the evolutionary selection of a new function.[35][36]
+
+To explain the observed specificity of enzymes, in 1894 Emil Fischer proposed that both the enzyme and the substrate possess specific complementary geometric shapes that fit exactly into one another.[37] This is often referred to as "the lock and key" model.[1]:8.3.2 This early model explains enzyme specificity, but fails to explain the stabilization of the transition state that enzymes achieve.[38]
+
+In 1958, Daniel Koshland suggested a modification to the lock and key model: since enzymes are rather flexible structures, the active site is continuously reshaped by interactions with the substrate as the substrate interacts with the enzyme.[39] As a result, the substrate does not simply bind to a rigid active site; the amino acid side-chains that make up the active site are molded into the precise positions that enable the enzyme to perform its catalytic function. In some cases, such as glycosidases, the substrate molecule also changes shape slightly as it enters the active site.[40] The active site continues to change until the substrate is completely bound, at which point the final shape and charge distribution is determined.[41]
+Induced fit may enhance the fidelity of molecular recognition in the presence of competition and noise via the conformational proofreading mechanism.[42]
+
+Enzymes can accelerate reactions in several ways, all of which lower the activation energy (ΔG‡, Gibbs free energy)[43]
+
+Enzymes may use several of these mechanisms simultaneously. For example, proteases such as trypsin perform covalent catalysis using a catalytic triad, stabilise charge build-up on the transition states using an oxyanion hole, complete hydrolysis using an oriented water substrate.[49]
+
+Enzymes are not rigid, static structures; instead they have complex internal dynamic motions – that is, movements of parts of the enzyme's structure such as individual amino acid residues, groups of residues forming a protein loop or unit of secondary structure, or even an entire protein domain. These motions give rise to a conformational ensemble of slightly different structures that interconvert with one another at equilibrium. Different states within this ensemble may be associated with different aspects of an enzyme's function. For example, different conformations of the enzyme dihydrofolate reductase are associated with the substrate binding, catalysis, cofactor release, and product release steps of the catalytic cycle,[50] consistent with catalytic resonance theory.
+
+Substrate presentation is a process where the enzyme is sequestered away from its substrate. Enzymes can be sequestered to the plasma membrane away from a substrate in the nucleus or cytosol. Or within the membrane, an enzyme can be sequestered into lipid rafts away from its substrate in the disordered region. When the enzyme is releases it mixes with its substrate. Alternatively, the enzyme can be sequestered near its substrate to activate the enzyme. For example, the enzyme can be soluble and upon activation bind to a lipid in the plasma membrane and then act upon molecules in the plasma membrane.
+
+Allosteric sites are pockets on the enzyme, distinct from the active site, that bind to molecules in the cellular environment. These molecules then cause a change in the conformation or dynamics of the enzyme that is transduced to the active site and thus affects the reaction rate of the enzyme.[51] In this way, allosteric interactions can either inhibit or activate enzymes. Allosteric interactions with metabolites upstream or downstream in an enzyme's metabolic pathway cause feedback regulation, altering the activity of the enzyme according to the flux through the rest of the pathway.[52]
+
+Some enzymes do not need additional components to show full activity. Others require non-protein molecules called cofactors to be bound for activity.[53] Cofactors can be either inorganic (e.g., metal ions and iron-sulfur clusters) or organic compounds (e.g., flavin and heme). These cofactors serve many purposes; for instance, metal ions can help in stabilizing nucleophilic species within the active site.[54] Organic cofactors can be either coenzymes, which are released from the enzyme's active site during the reaction, or prosthetic groups, which are tightly bound to an enzyme. Organic prosthetic groups can be covalently bound (e.g., biotin in enzymes such as pyruvate carboxylase).[55]
+
+An example of an enzyme that contains a cofactor is carbonic anhydrase, which uses a zinc cofactor bound as part of its active site.[56] These tightly bound ions or molecules are usually found in the active site and are involved in catalysis.[1]:8.1.1 For example, flavin and heme cofactors are often involved in redox reactions.[1]:17
+
+Enzymes that require a cofactor but do not have one bound are called apoenzymes or apoproteins. An enzyme together with the cofactor(s) required for activity is called a holoenzyme (or haloenzyme). The term holoenzyme can also be applied to enzymes that contain multiple protein subunits, such as the DNA polymerases; here the holoenzyme is the complete complex containing all the subunits needed for activity.[1]:8.1.1
+
+Coenzymes are small organic molecules that can be loosely or tightly bound to an enzyme. Coenzymes transport chemical groups from one enzyme to another.[57] Examples include NADH, NADPH and adenosine triphosphate (ATP). Some coenzymes, such as flavin mononucleotide (FMN), flavin adenine dinucleotide (FAD), thiamine pyrophosphate (TPP), and tetrahydrofolate (THF), are derived from vitamins. These coenzymes cannot be synthesized by the body de novo and closely related compounds (vitamins) must be acquired from the diet. The chemical groups carried include:
+
+Since coenzymes are chemically changed as a consequence of enzyme action, it is useful to consider coenzymes to be a special class of substrates, or second substrates, which are common to many different enzymes. For example, about 1000 enzymes are known to use the coenzyme NADH.[58]
+
+Coenzymes are usually continuously regenerated and their concentrations maintained at a steady level inside the cell. For example, NADPH is regenerated through the pentose phosphate pathway and S-adenosylmethionine by methionine adenosyltransferase. This continuous regeneration means that small amounts of coenzymes can be used very intensively. For example, the human body turns over its own weight in ATP each day.[59]
+
+As with all catalysts, enzymes do not alter the position of the chemical equilibrium of the reaction. In the presence of an enzyme, the reaction runs in the same direction as it would without the enzyme, just more quickly.[1]:8.2.3 For example, carbonic anhydrase catalyzes its reaction in either direction depending on the concentration of its reactants:[60]
+
+
+
+
+
+
+
+
+
+(1)
+
+
+
+
+
+
+
+
+
+(2)
+
+The rate of a reaction is dependent on the activation energy needed to form the transition state which then decays into products. Enzymes increase reaction rates by lowering the energy of the transition state. First, binding forms a low energy enzyme-substrate complex (ES). Second, the enzyme stabilises the transition state such that it requires less energy to achieve compared to the uncatalyzed reaction (ES‡). Finally the enzyme-product complex (EP) dissociates to release the products.[1]:8.3
+
+Enzymes can couple two or more reactions, so that a thermodynamically favorable reaction can be used to "drive" a thermodynamically unfavourable one so that the combined energy of the products is lower than the substrates. For example, the hydrolysis of ATP is often used to drive other chemical reactions.[61]
+
+Enzyme kinetics is the investigation of how enzymes bind substrates and turn them into products.[62] The rate data used in kinetic analyses are commonly obtained from enzyme assays. In 1913 Leonor Michaelis and Maud Leonora Menten proposed a quantitative theory of enzyme kinetics, which is referred to as Michaelis–Menten kinetics.[63] The major contribution of Michaelis and Menten was to think of enzyme reactions in two stages. In the first, the substrate binds reversibly to the enzyme, forming the enzyme-substrate complex. This is sometimes called the Michaelis–Menten complex in their honor. The enzyme then catalyzes the chemical step in the reaction and releases the product. This work was further developed by G. E. Briggs and J. B. S. Haldane, who derived kinetic equations that are still widely used today.[64]
+
+Enzyme rates depend on solution conditions and substrate concentration. To find the maximum speed of an enzymatic reaction, the substrate concentration is increased until a constant rate of product formation is seen. This is shown in the saturation curve on the right. Saturation happens because, as substrate concentration increases, more and more of the free enzyme is converted into the substrate-bound ES complex. At the maximum reaction rate (Vmax) of the enzyme, all the enzyme active sites are bound to substrate, and the amount of ES complex is the same as the total amount of enzyme.[1]:8.4
+
+Vmax is only one of several important kinetic parameters. The amount of substrate needed to achieve a given rate of reaction is also important. This is given by the Michaelis–Menten constant (Km), which is the substrate concentration required for an enzyme to reach one-half its maximum reaction rate; generally, each enzyme has a characteristic KM for a given substrate. Another useful constant is kcat, also called the turnover number, which is the number of substrate molecules handled by one active site per second.[1]:8.4
+
+The efficiency of an enzyme can be expressed in terms of kcat/Km. This is also called the specificity constant and incorporates the rate constants for all steps in the reaction up to and including the first irreversible step. Because the specificity constant reflects both affinity and catalytic ability, it is useful for comparing different enzymes against each other, or the same enzyme with different substrates. The theoretical maximum for the specificity constant is called the diffusion limit and is about 108 to 109 (M−1 s−1). At this point every collision of the enzyme with its substrate will result in catalysis, and the rate of product formation is not limited by the reaction rate but by the diffusion rate. Enzymes with this property are called catalytically perfect or kinetically perfect. Example of such enzymes are triose-phosphate isomerase, carbonic anhydrase, acetylcholinesterase, catalase, fumarase, β-lactamase, and superoxide dismutase.[1]:8.4.2 The turnover of such enzymes can reach several million reactions per second.[1]:9.2 But most enzymes are far from perfect: the average values of 
+
+
+
+
+k
+
+
+c
+a
+t
+
+
+
+
+/
+
+
+K
+
+
+m
+
+
+
+
+
+{\displaystyle k_{\rm {cat}}/K_{\rm {m}}}
+
+ and 
+
+
+
+
+k
+
+
+c
+a
+t
+
+
+
+
+
+{\displaystyle k_{\rm {cat}}}
+
+ are about 
+
+
+
+
+10
+
+5
+
+
+
+
+
+s
+
+
+
+−
+1
+
+
+
+
+
+M
+
+
+
+−
+1
+
+
+
+
+{\displaystyle 10^{5}{\rm {s}}^{-1}{\rm {M}}^{-1}}
+
+ and 
+
+
+
+10
+
+
+
+s
+
+
+
+−
+1
+
+
+
+
+{\displaystyle 10{\rm {s}}^{-1}}
+
+, respectively.[65]
+
+Michaelis–Menten kinetics relies on the law of mass action, which is derived from the assumptions of free diffusion and thermodynamically driven random collision. Many biochemical or cellular processes deviate significantly from these conditions, because of macromolecular crowding and constrained molecular movement.[66] More recent, complex extensions of the model attempt to correct for these effects.[67]
+
+Enzyme reaction rates can be decreased by various types of enzyme inhibitors.[69]:73–74
+
+A competitive inhibitor and substrate cannot bind to the enzyme at the same time.[70] Often competitive inhibitors strongly resemble the real substrate of the enzyme. For example, the drug methotrexate is a competitive inhibitor of the enzyme dihydrofolate reductase, which catalyzes the reduction of dihydrofolate to tetrahydrofolate.[68] The similarity between the structures of dihydrofolate and this drug are shown in the accompanying figure. This type of inhibition can be overcome with high substrate concentration. In some cases, the inhibitor can bind to a site other than the binding-site of the usual substrate and exert an allosteric effect to change the shape of the usual binding-site.[71]
+
+A non-competitive inhibitor binds to a site other than where the substrate binds. The substrate still binds with its usual affinity and hence Km remains the same. However the inhibitor reduces the catalytic efficiency of the enzyme so that Vmax is reduced. In contrast to competitive inhibition, non-competitive inhibition cannot be overcome with high substrate concentration.[69]:76–78
+
+An uncompetitive inhibitor cannot bind to the free enzyme, only to the enzyme-substrate complex; hence, these types of inhibitors are most effective at high substrate concentration. In the presence of the inhibitor, the enzyme-substrate complex is inactive.[69]:78 This type of inhibition is rare.[72]
+
+A mixed inhibitor binds to an allosteric site and the binding of the substrate and the inhibitor affect each other. The enzyme's function is reduced but not eliminated when bound to the inhibitor. This type of inhibitor does not follow the Michaelis–Menten equation.[69]:76–78
+
+An irreversible inhibitor permanently inactivates the enzyme, usually by forming a covalent bond to the protein.[73] Penicillin[74] and aspirin[75] are common drugs that act in this manner.
+
+In many organisms, inhibitors may act as part of a feedback mechanism. If an enzyme produces too much of one substance in the organism, that substance may act as an inhibitor for the enzyme at the beginning of the pathway that produces it, causing production of the substance to slow down or stop when there is sufficient amount. This is a form of negative feedback. Major metabolic pathways such as the citric acid cycle make use of this mechanism.[1]:17.2.2
+
+Since inhibitors modulate the function of enzymes they are often used as drugs. Many such drugs are reversible competitive inhibitors that resemble the enzyme's native substrate, similar to methotrexate above; other well-known examples include statins used to treat high cholesterol,[76] and protease inhibitors used to treat retroviral infections such as HIV.[77] A common example of an irreversible inhibitor that is used as a drug is aspirin, which inhibits the COX-1 and COX-2 enzymes that produce the inflammation messenger prostaglandin.[75] Other enzyme inhibitors are poisons. For example, the poison cyanide is an irreversible enzyme inhibitor that combines with the copper and iron in the active site of the enzyme cytochrome c oxidase and blocks cellular respiration.[78]
+
+As enzymes are made up of proteins, their actions are sensitive to change in many physio chemical factors such as pH, temperature, substrate concentration, etc.
+
+The following table shows pH optima for various enzymes.[79]
+
+Enzymes serve a wide variety of functions inside living organisms. They are indispensable for signal transduction and cell regulation, often via kinases and phosphatases.[80] They also generate movement, with myosin hydrolyzing ATP to generate muscle contraction, and also transport cargo around the cell as part of the cytoskeleton.[81] Other ATPases in the cell membrane are ion pumps involved in active transport. Enzymes are also involved in more exotic functions, such as luciferase generating light in fireflies.[82] Viruses can also contain enzymes for infecting cells, such as the HIV integrase and reverse transcriptase, or for viral release from cells, like the influenza virus neuraminidase.[83]
+
+An important function of enzymes is in the digestive systems of animals. Enzymes such as amylases and proteases break down large molecules (starch or proteins, respectively) into smaller ones, so they can be absorbed by the intestines. Starch molecules, for example, are too large to be absorbed from the intestine, but enzymes hydrolyze the starch chains into smaller molecules such as maltose and eventually glucose, which can then be absorbed. Different enzymes digest different food substances. In ruminants, which have herbivorous diets, microorganisms in the gut produce another enzyme, cellulase, to break down the cellulose cell walls of plant fiber.[84]
+
+Several enzymes can work together in a specific order, creating metabolic pathways.[1]:30.1 In a metabolic pathway, one enzyme takes the product of another enzyme as a substrate. After the catalytic reaction, the product is then passed on to another enzyme. Sometimes more than one enzyme can catalyze the same reaction in parallel; this can allow more complex regulation: with, for example, a low constant activity provided by one enzyme but an inducible high activity from a second enzyme.[85]
+
+Enzymes determine what steps occur in these pathways. Without enzymes, metabolism would neither progress through the same steps and could not be regulated to serve the needs of the cell. Most central metabolic pathways are regulated at a few key steps, typically through enzymes whose activity involves the hydrolysis of ATP. Because this reaction releases so much energy, other reactions that are thermodynamically unfavorable can be coupled to ATP hydrolysis, driving the overall series of linked metabolic reactions.[1]:30.1
+
+There are five main ways that enzyme activity is controlled in the cell.[1]:30.1.1
+
+Enzymes can be either activated or inhibited by other molecules. For example, the end product(s) of a metabolic pathway are often inhibitors for one of the first enzymes of the pathway (usually the first irreversible step, called committed step), thus regulating the amount of end product made by the pathways. Such a regulatory mechanism is called a negative feedback mechanism, because the amount of the end product produced is regulated by its own concentration.[86]:141–48 Negative feedback mechanism can effectively adjust the rate of synthesis of intermediate metabolites according to the demands of the cells. This helps with effective allocations of materials and energy economy, and it prevents the excess manufacture of end products. Like other homeostatic devices, the control of enzymatic action helps to maintain a stable internal environment in living organisms.[86]:141
+
+Examples of post-translational modification include phosphorylation, myristoylation and glycosylation.[86]:149–69 For example, in the response to insulin, the phosphorylation of multiple enzymes, including glycogen synthase, helps control the synthesis or degradation of glycogen and allows the cell to respond to changes in blood sugar.[87] Another example of post-translational modification is the cleavage of the polypeptide chain. Chymotrypsin, a digestive protease, is produced in inactive form as chymotrypsinogen in the pancreas and transported in this form to the stomach where it is activated. This stops the enzyme from digesting the pancreas or other tissues before it enters the gut. This type of inactive precursor to an enzyme is known as a zymogen[86]:149–53 or proenzyme.
+
+Enzyme production (transcription and translation of enzyme genes) can be enhanced or diminished by a cell in response to changes in the cell's environment. This form of gene regulation is called enzyme induction. For example, bacteria may become resistant to antibiotics such as penicillin because enzymes called beta-lactamases are induced that hydrolyse the crucial beta-lactam ring within the penicillin molecule.[88] Another example comes from enzymes in the liver called cytochrome P450 oxidases, which are important in drug metabolism. Induction or inhibition of these enzymes can cause drug interactions.[89] Enzyme levels can also be regulated by changing the rate of enzyme degradation.[1]:30.1.1 The opposite of enzyme induction is enzyme repression.
+
+Enzymes can be compartmentalized, with different metabolic pathways occurring in different cellular compartments. For example, fatty acids are synthesized by one set of enzymes in the cytosol, endoplasmic reticulum and Golgi and used by a different set of enzymes as a source of energy in the mitochondrion, through β-oxidation.[90] In addition, trafficking of the enzyme to different compartments may change the degree of protonation (e.g., the neutral cytoplasm and the acidic lysosome) or oxidative state (e.g., oxidizing periplasm or reducing cytoplasm) which in turn affects enzyme activity.[91] In contrast to partitioning into membrane bound organelles, enzyme subcellular localisation may also be altered through polymerisation of enzymes into macromolecular cytoplasmic filaments.[92][93]
+
+In multicellular eukaryotes, cells in different organs and tissues have different patterns of gene expression and therefore have different sets of enzymes (known as isozymes) available for metabolic reactions. This provides a mechanism for regulating the overall metabolism of the organism. For example, hexokinase, the first enzyme in the glycolysis pathway, has a specialized form called glucokinase expressed in the liver and pancreas that has a lower affinity for glucose yet is more sensitive to glucose concentration.[94] This enzyme is involved in sensing blood sugar and regulating insulin production.[95]
+
+Since the tight control of enzyme activity is essential for homeostasis, any malfunction (mutation, overproduction, underproduction or deletion) of a single critical enzyme can lead to a genetic disease. The malfunction of just one type of enzyme out of the thousands of types present in the human body can be fatal. An example of a fatal genetic disease due to enzyme insufficiency is Tay–Sachs disease, in which patients lack the enzyme hexosaminidase.[96][97]
+
+One example of enzyme deficiency is the most common type of phenylketonuria. Many different single amino acid mutations in the enzyme phenylalanine hydroxylase, which catalyzes the first step in the degradation of phenylalanine, result in build-up of phenylalanine and related products. Some mutations are in the active site, directly disrupting binding and catalysis, but many are far from the active site and reduce activity by destabilising the protein structure, or affecting correct oligomerisation.[98][99] This can lead to intellectual disability if the disease is untreated.[100] Another example is pseudocholinesterase deficiency, in which the body's ability to break down choline ester drugs is impaired.[101] 
+Oral administration of enzymes can be used to treat some functional enzyme deficiencies, such as pancreatic insufficiency[102] and lactose intolerance.[103]
+
+Another way enzyme malfunctions can cause disease comes from germline mutations in genes coding for DNA repair enzymes. Defects in these enzymes cause cancer because cells are less able to repair mutations in their genomes. This causes a slow accumulation of mutations and results in the development of cancers. An example of such a hereditary cancer syndrome is xeroderma pigmentosum, which causes the development of skin cancers in response to even minimal exposure to ultraviolet light.[104][105]
+
+Similar to any other protein, enzymes change over time through mutations and sequence divergence. Given their central role in metabolism, enzyme evolution plays a critical role in adaptation. A key question is therefore whether and how enzymes can change their enzymatic activities alongside. It is generally accepted that many new enzyme activities have evolved through gene duplication and mutation of the duplicate copies although evolution can also happen without duplication. One example of an enzyme that has changed its activity is the ancestor of methionyl amino peptidase (MAP) and creatine amidinohydrolase (creatinase) which are clearly homologous but catalyze very different reactions (MAP removes the amino-terminal methionine in new proteins while creatinase hydrolyses creatine to sarcosine and urea). In addition, MAP is metal-ion dependent while creatinase is not, hence this property was also lost over time.[106] Small changes of enzymatic activity are extremely common among enzymes. In particular, substrate binding specificity (see above) can easily and quickly change with single amino acid changes in their substrate binding pockets. This is frequently seen in the main enzyme classes such as kinases.[107]
+
+Artificial (in vitro) evolution is now commonly used to modify enzyme activity or specificity for industrial applications (see below).
+
+Enzymes are used in the chemical industry and other industrial applications when extremely specific catalysts are required. Enzymes in general are limited in the number of reactions they have evolved to catalyze and also by their lack of stability in organic solvents and at high temperatures. As a consequence, protein engineering is an active area of research and involves attempts to create new enzymes with novel properties, either through rational design or in vitro evolution.[108][109] These efforts have begun to be successful, and a few enzymes have now been designed "from scratch" to catalyze reactions that do not occur in nature.[110]
+
+General
+
+Etymology and history
+
+
+
+Enzyme structure and mechanism
+
+Kinetics and inhibition
+
+
+

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+Enzymes /ˈɛnzaɪmz/ are proteins that act as biological catalysts (biocatalysts). Catalysts accelerate chemical reactions. The molecules upon which enzymes may act are called substrates, and the enzyme converts the substrates into different molecules known as products. Almost all metabolic processes in the cell need enzyme catalysis in order to occur at rates fast enough to sustain life.[1]:8.1 Metabolic pathways depend upon enzymes to catalyze individual steps. The study of enzymes is called enzymology and a new field of pseudoenzyme analysis has recently grown up, recognising that during evolution, some enzymes have lost the ability to carry out biological catalysis, which is often reflected in their amino acid sequences and unusual 'pseudocatalytic' properties.[2][3]
+
+Enzymes are known to catalyze more than 5,000 biochemical reaction types.[4] Other biocatalysts are catalytic RNA molecules, called ribozymes. Enzymes' specificity comes from their unique three-dimensional structures.
+
+Like all catalysts, enzymes increase the reaction rate by lowering its activation energy. Some enzymes can make their conversion of substrate to product occur many millions of times faster. An extreme example is orotidine 5'-phosphate decarboxylase, which allows a reaction that would otherwise take millions of years to occur in milliseconds.[5][6] Chemically, enzymes are like any catalyst and are not consumed in chemical reactions, nor do they alter the equilibrium of a reaction. Enzymes differ from most other catalysts by being much more specific. Enzyme activity can be affected by other molecules: inhibitors are molecules that decrease enzyme activity, and activators are molecules that increase activity. Many therapeutic drugs and poisons are enzyme inhibitors. An enzyme's activity decreases markedly outside its optimal temperature and pH, and many enzymes are (permanently) denatured when exposed to excessive heat, losing their structure and catalytic properties.
+
+Some enzymes are used commercially, for example, in the synthesis of antibiotics. Some household products use enzymes to speed up chemical reactions: enzymes in biological washing powders break down protein, starch or fat stains on clothes, and enzymes in meat tenderizer break down proteins into smaller molecules, making the meat easier to chew.
+
+By the late 17th and early 18th centuries, the digestion of meat by stomach secretions[7] and the conversion of starch to sugars by plant extracts and saliva were known but the mechanisms by which these occurred had not been identified.[8]
+
+French chemist Anselme Payen was the first to discover an enzyme, diastase, in 1833.[9] A few decades later, when studying the fermentation of sugar to alcohol by yeast, Louis Pasteur concluded that this fermentation was caused by a vital force contained within the yeast cells called "ferments", which were thought to function only within living organisms. He wrote that "alcoholic fermentation is an act correlated with the life and organization of the yeast cells, not with the death or putrefaction of the cells."[10]
+
+In 1877, German physiologist Wilhelm Kühne (1837–1900) first used the term enzyme, which comes from Greek ἔνζυμον, "leavened" or "in yeast", to describe this process.[11] The word enzyme was used later to refer to nonliving substances such as pepsin, and the word ferment was used to refer to chemical activity produced by living organisms.[12]
+
+Eduard Buchner submitted his first paper on the study of yeast extracts in 1897. In a series of experiments at the University of Berlin, he found that sugar was fermented by yeast extracts even when there were no living yeast cells in the mixture.[13] He named the enzyme that brought about the fermentation of sucrose "zymase".[14] In 1907, he received the Nobel Prize in Chemistry for "his discovery of cell-free fermentation". Following Buchner's example, enzymes are usually named according to the reaction they carry out: the suffix -ase is combined with the name of the substrate (e.g., lactase is the enzyme that cleaves lactose) or to the type of reaction (e.g., DNA polymerase forms DNA polymers).[15]
+
+The biochemical identity of enzymes was still unknown in the early 1900s. Many scientists observed that enzymatic activity was associated with proteins, but others (such as Nobel laureate Richard Willstätter) argued that proteins were merely carriers for the true enzymes and that proteins per se were incapable of catalysis.[16] In 1926, James B. Sumner showed that the enzyme urease was a pure protein and crystallized it; he did likewise for the enzyme catalase in 1937. The conclusion that pure proteins can be enzymes was definitively demonstrated by John Howard Northrop and Wendell Meredith Stanley, who worked on the digestive enzymes pepsin (1930), trypsin and chymotrypsin. These three scientists were awarded the 1946 Nobel Prize in Chemistry.[17]
+
+The discovery that enzymes could be crystallized eventually allowed their structures to be solved by x-ray crystallography. This was first done for lysozyme, an enzyme found in tears, saliva and egg whites that digests the coating of some bacteria; the structure was solved by a group led by David Chilton Phillips and published in 1965.[18] This high-resolution structure of lysozyme marked the beginning of the field of structural biology and the effort to understand how enzymes work at an atomic level of detail.[19]
+
+An enzyme's name is often derived from its substrate or the chemical reaction it catalyzes, with the word ending in -ase.[1]:8.1.3 Examples are lactase, alcohol dehydrogenase and DNA polymerase. Different enzymes that catalyze the same chemical reaction are called isozymes.[1]:10.3
+
+The International Union of Biochemistry and Molecular Biology have developed a nomenclature for enzymes, the EC numbers; each enzyme is described by a sequence of four numbers preceded by "EC", which stands for "Enzyme Commission". The first number broadly classifies the enzyme based on its mechanism.[20]
+
+The top-level classification is:
+
+These sections are subdivided by other features such as the substrate, products, and chemical mechanism. An enzyme is fully specified by four numerical designations. For example, hexokinase (EC 2.7.1.1) is a transferase (EC 2) that adds a phosphate group (EC 2.7) to a hexose sugar, a molecule containing an alcohol group (EC 2.7.1).[21]
+
+Enzymes are generally globular proteins, acting alone or in larger complexes. The sequence of the amino acids specifies the structure which in turn determines the catalytic activity of the enzyme.[22] Although structure determines function, a novel enzymatic activity cannot yet be predicted from structure alone.[23] Enzyme structures unfold (denature) when heated or exposed to chemical denaturants and this disruption to the structure typically causes a loss of activity.[24] Enzyme denaturation is normally linked to temperatures above a species' normal level; as a result, enzymes from bacteria living in volcanic environments such as hot springs are prized by industrial users for their ability to function at high temperatures, allowing enzyme-catalysed reactions to be operated at a very high rate.
+
+Enzymes are usually much larger than their substrates. Sizes range from just 62 amino acid residues, for the monomer of 4-oxalocrotonate tautomerase,[25] to over 2,500 residues in the animal fatty acid synthase.[26] Only a small portion of their structure (around 2–4 amino acids) is directly involved in catalysis: the catalytic site.[27] This catalytic site is located next to one or more binding sites where residues orient the substrates. The catalytic site and binding site together compose the enzyme's active site. The remaining majority of the enzyme structure serves to maintain the precise orientation and dynamics of the active site.[28]
+
+In some enzymes, no amino acids are directly involved in catalysis; instead, the enzyme contains sites to bind and orient catalytic cofactors.[28] Enzyme structures may also contain allosteric sites where the binding of a small molecule causes a conformational change that increases or decreases activity.[29]
+
+A small number of RNA-based biological catalysts called ribozymes exist, which again can act alone or in complex with proteins. The most common of these is the ribosome which is a complex of protein and catalytic RNA components.[1]:2.2
+
+Enzymes must bind their substrates before they can catalyse any chemical reaction. Enzymes are usually very specific as to what substrates they bind and then the chemical reaction catalysed. Specificity is achieved by binding pockets with complementary shape, charge and hydrophilic/hydrophobic characteristics to the substrates. Enzymes can therefore distinguish between very similar substrate molecules to be chemoselective, regioselective and stereospecific.[30]
+
+Some of the enzymes showing the highest specificity and accuracy are involved in the copying and expression of the genome. Some of these enzymes have "proof-reading" mechanisms. Here, an enzyme such as DNA polymerase catalyzes a reaction in a first step and then checks that the product is correct in a second step.[31] This two-step process results in average error rates of less than 1 error in 100 million reactions in high-fidelity mammalian polymerases.[1]:5.3.1 Similar proofreading mechanisms are also found in RNA polymerase,[32] aminoacyl tRNA synthetases[33] and ribosomes.[34]
+
+Conversely, some enzymes display enzyme promiscuity, having broad specificity and acting on a range of different physiologically relevant substrates. Many enzymes possess small side activities which arose fortuitously (i.e. neutrally), which may be the starting point for the evolutionary selection of a new function.[35][36]
+
+To explain the observed specificity of enzymes, in 1894 Emil Fischer proposed that both the enzyme and the substrate possess specific complementary geometric shapes that fit exactly into one another.[37] This is often referred to as "the lock and key" model.[1]:8.3.2 This early model explains enzyme specificity, but fails to explain the stabilization of the transition state that enzymes achieve.[38]
+
+In 1958, Daniel Koshland suggested a modification to the lock and key model: since enzymes are rather flexible structures, the active site is continuously reshaped by interactions with the substrate as the substrate interacts with the enzyme.[39] As a result, the substrate does not simply bind to a rigid active site; the amino acid side-chains that make up the active site are molded into the precise positions that enable the enzyme to perform its catalytic function. In some cases, such as glycosidases, the substrate molecule also changes shape slightly as it enters the active site.[40] The active site continues to change until the substrate is completely bound, at which point the final shape and charge distribution is determined.[41]
+Induced fit may enhance the fidelity of molecular recognition in the presence of competition and noise via the conformational proofreading mechanism.[42]
+
+Enzymes can accelerate reactions in several ways, all of which lower the activation energy (ΔG‡, Gibbs free energy)[43]
+
+Enzymes may use several of these mechanisms simultaneously. For example, proteases such as trypsin perform covalent catalysis using a catalytic triad, stabilise charge build-up on the transition states using an oxyanion hole, complete hydrolysis using an oriented water substrate.[49]
+
+Enzymes are not rigid, static structures; instead they have complex internal dynamic motions – that is, movements of parts of the enzyme's structure such as individual amino acid residues, groups of residues forming a protein loop or unit of secondary structure, or even an entire protein domain. These motions give rise to a conformational ensemble of slightly different structures that interconvert with one another at equilibrium. Different states within this ensemble may be associated with different aspects of an enzyme's function. For example, different conformations of the enzyme dihydrofolate reductase are associated with the substrate binding, catalysis, cofactor release, and product release steps of the catalytic cycle,[50] consistent with catalytic resonance theory.
+
+Substrate presentation is a process where the enzyme is sequestered away from its substrate. Enzymes can be sequestered to the plasma membrane away from a substrate in the nucleus or cytosol. Or within the membrane, an enzyme can be sequestered into lipid rafts away from its substrate in the disordered region. When the enzyme is releases it mixes with its substrate. Alternatively, the enzyme can be sequestered near its substrate to activate the enzyme. For example, the enzyme can be soluble and upon activation bind to a lipid in the plasma membrane and then act upon molecules in the plasma membrane.
+
+Allosteric sites are pockets on the enzyme, distinct from the active site, that bind to molecules in the cellular environment. These molecules then cause a change in the conformation or dynamics of the enzyme that is transduced to the active site and thus affects the reaction rate of the enzyme.[51] In this way, allosteric interactions can either inhibit or activate enzymes. Allosteric interactions with metabolites upstream or downstream in an enzyme's metabolic pathway cause feedback regulation, altering the activity of the enzyme according to the flux through the rest of the pathway.[52]
+
+Some enzymes do not need additional components to show full activity. Others require non-protein molecules called cofactors to be bound for activity.[53] Cofactors can be either inorganic (e.g., metal ions and iron-sulfur clusters) or organic compounds (e.g., flavin and heme). These cofactors serve many purposes; for instance, metal ions can help in stabilizing nucleophilic species within the active site.[54] Organic cofactors can be either coenzymes, which are released from the enzyme's active site during the reaction, or prosthetic groups, which are tightly bound to an enzyme. Organic prosthetic groups can be covalently bound (e.g., biotin in enzymes such as pyruvate carboxylase).[55]
+
+An example of an enzyme that contains a cofactor is carbonic anhydrase, which uses a zinc cofactor bound as part of its active site.[56] These tightly bound ions or molecules are usually found in the active site and are involved in catalysis.[1]:8.1.1 For example, flavin and heme cofactors are often involved in redox reactions.[1]:17
+
+Enzymes that require a cofactor but do not have one bound are called apoenzymes or apoproteins. An enzyme together with the cofactor(s) required for activity is called a holoenzyme (or haloenzyme). The term holoenzyme can also be applied to enzymes that contain multiple protein subunits, such as the DNA polymerases; here the holoenzyme is the complete complex containing all the subunits needed for activity.[1]:8.1.1
+
+Coenzymes are small organic molecules that can be loosely or tightly bound to an enzyme. Coenzymes transport chemical groups from one enzyme to another.[57] Examples include NADH, NADPH and adenosine triphosphate (ATP). Some coenzymes, such as flavin mononucleotide (FMN), flavin adenine dinucleotide (FAD), thiamine pyrophosphate (TPP), and tetrahydrofolate (THF), are derived from vitamins. These coenzymes cannot be synthesized by the body de novo and closely related compounds (vitamins) must be acquired from the diet. The chemical groups carried include:
+
+Since coenzymes are chemically changed as a consequence of enzyme action, it is useful to consider coenzymes to be a special class of substrates, or second substrates, which are common to many different enzymes. For example, about 1000 enzymes are known to use the coenzyme NADH.[58]
+
+Coenzymes are usually continuously regenerated and their concentrations maintained at a steady level inside the cell. For example, NADPH is regenerated through the pentose phosphate pathway and S-adenosylmethionine by methionine adenosyltransferase. This continuous regeneration means that small amounts of coenzymes can be used very intensively. For example, the human body turns over its own weight in ATP each day.[59]
+
+As with all catalysts, enzymes do not alter the position of the chemical equilibrium of the reaction. In the presence of an enzyme, the reaction runs in the same direction as it would without the enzyme, just more quickly.[1]:8.2.3 For example, carbonic anhydrase catalyzes its reaction in either direction depending on the concentration of its reactants:[60]
+
+
+
+
+
+
+
+
+
+(1)
+
+
+
+
+
+
+
+
+
+(2)
+
+The rate of a reaction is dependent on the activation energy needed to form the transition state which then decays into products. Enzymes increase reaction rates by lowering the energy of the transition state. First, binding forms a low energy enzyme-substrate complex (ES). Second, the enzyme stabilises the transition state such that it requires less energy to achieve compared to the uncatalyzed reaction (ES‡). Finally the enzyme-product complex (EP) dissociates to release the products.[1]:8.3
+
+Enzymes can couple two or more reactions, so that a thermodynamically favorable reaction can be used to "drive" a thermodynamically unfavourable one so that the combined energy of the products is lower than the substrates. For example, the hydrolysis of ATP is often used to drive other chemical reactions.[61]
+
+Enzyme kinetics is the investigation of how enzymes bind substrates and turn them into products.[62] The rate data used in kinetic analyses are commonly obtained from enzyme assays. In 1913 Leonor Michaelis and Maud Leonora Menten proposed a quantitative theory of enzyme kinetics, which is referred to as Michaelis–Menten kinetics.[63] The major contribution of Michaelis and Menten was to think of enzyme reactions in two stages. In the first, the substrate binds reversibly to the enzyme, forming the enzyme-substrate complex. This is sometimes called the Michaelis–Menten complex in their honor. The enzyme then catalyzes the chemical step in the reaction and releases the product. This work was further developed by G. E. Briggs and J. B. S. Haldane, who derived kinetic equations that are still widely used today.[64]
+
+Enzyme rates depend on solution conditions and substrate concentration. To find the maximum speed of an enzymatic reaction, the substrate concentration is increased until a constant rate of product formation is seen. This is shown in the saturation curve on the right. Saturation happens because, as substrate concentration increases, more and more of the free enzyme is converted into the substrate-bound ES complex. At the maximum reaction rate (Vmax) of the enzyme, all the enzyme active sites are bound to substrate, and the amount of ES complex is the same as the total amount of enzyme.[1]:8.4
+
+Vmax is only one of several important kinetic parameters. The amount of substrate needed to achieve a given rate of reaction is also important. This is given by the Michaelis–Menten constant (Km), which is the substrate concentration required for an enzyme to reach one-half its maximum reaction rate; generally, each enzyme has a characteristic KM for a given substrate. Another useful constant is kcat, also called the turnover number, which is the number of substrate molecules handled by one active site per second.[1]:8.4
+
+The efficiency of an enzyme can be expressed in terms of kcat/Km. This is also called the specificity constant and incorporates the rate constants for all steps in the reaction up to and including the first irreversible step. Because the specificity constant reflects both affinity and catalytic ability, it is useful for comparing different enzymes against each other, or the same enzyme with different substrates. The theoretical maximum for the specificity constant is called the diffusion limit and is about 108 to 109 (M−1 s−1). At this point every collision of the enzyme with its substrate will result in catalysis, and the rate of product formation is not limited by the reaction rate but by the diffusion rate. Enzymes with this property are called catalytically perfect or kinetically perfect. Example of such enzymes are triose-phosphate isomerase, carbonic anhydrase, acetylcholinesterase, catalase, fumarase, β-lactamase, and superoxide dismutase.[1]:8.4.2 The turnover of such enzymes can reach several million reactions per second.[1]:9.2 But most enzymes are far from perfect: the average values of 
+
+
+
+
+k
+
+
+c
+a
+t
+
+
+
+
+/
+
+
+K
+
+
+m
+
+
+
+
+
+{\displaystyle k_{\rm {cat}}/K_{\rm {m}}}
+
+ and 
+
+
+
+
+k
+
+
+c
+a
+t
+
+
+
+
+
+{\displaystyle k_{\rm {cat}}}
+
+ are about 
+
+
+
+
+10
+
+5
+
+
+
+
+
+s
+
+
+
+−
+1
+
+
+
+
+
+M
+
+
+
+−
+1
+
+
+
+
+{\displaystyle 10^{5}{\rm {s}}^{-1}{\rm {M}}^{-1}}
+
+ and 
+
+
+
+10
+
+
+
+s
+
+
+
+−
+1
+
+
+
+
+{\displaystyle 10{\rm {s}}^{-1}}
+
+, respectively.[65]
+
+Michaelis–Menten kinetics relies on the law of mass action, which is derived from the assumptions of free diffusion and thermodynamically driven random collision. Many biochemical or cellular processes deviate significantly from these conditions, because of macromolecular crowding and constrained molecular movement.[66] More recent, complex extensions of the model attempt to correct for these effects.[67]
+
+Enzyme reaction rates can be decreased by various types of enzyme inhibitors.[69]:73–74
+
+A competitive inhibitor and substrate cannot bind to the enzyme at the same time.[70] Often competitive inhibitors strongly resemble the real substrate of the enzyme. For example, the drug methotrexate is a competitive inhibitor of the enzyme dihydrofolate reductase, which catalyzes the reduction of dihydrofolate to tetrahydrofolate.[68] The similarity between the structures of dihydrofolate and this drug are shown in the accompanying figure. This type of inhibition can be overcome with high substrate concentration. In some cases, the inhibitor can bind to a site other than the binding-site of the usual substrate and exert an allosteric effect to change the shape of the usual binding-site.[71]
+
+A non-competitive inhibitor binds to a site other than where the substrate binds. The substrate still binds with its usual affinity and hence Km remains the same. However the inhibitor reduces the catalytic efficiency of the enzyme so that Vmax is reduced. In contrast to competitive inhibition, non-competitive inhibition cannot be overcome with high substrate concentration.[69]:76–78
+
+An uncompetitive inhibitor cannot bind to the free enzyme, only to the enzyme-substrate complex; hence, these types of inhibitors are most effective at high substrate concentration. In the presence of the inhibitor, the enzyme-substrate complex is inactive.[69]:78 This type of inhibition is rare.[72]
+
+A mixed inhibitor binds to an allosteric site and the binding of the substrate and the inhibitor affect each other. The enzyme's function is reduced but not eliminated when bound to the inhibitor. This type of inhibitor does not follow the Michaelis–Menten equation.[69]:76–78
+
+An irreversible inhibitor permanently inactivates the enzyme, usually by forming a covalent bond to the protein.[73] Penicillin[74] and aspirin[75] are common drugs that act in this manner.
+
+In many organisms, inhibitors may act as part of a feedback mechanism. If an enzyme produces too much of one substance in the organism, that substance may act as an inhibitor for the enzyme at the beginning of the pathway that produces it, causing production of the substance to slow down or stop when there is sufficient amount. This is a form of negative feedback. Major metabolic pathways such as the citric acid cycle make use of this mechanism.[1]:17.2.2
+
+Since inhibitors modulate the function of enzymes they are often used as drugs. Many such drugs are reversible competitive inhibitors that resemble the enzyme's native substrate, similar to methotrexate above; other well-known examples include statins used to treat high cholesterol,[76] and protease inhibitors used to treat retroviral infections such as HIV.[77] A common example of an irreversible inhibitor that is used as a drug is aspirin, which inhibits the COX-1 and COX-2 enzymes that produce the inflammation messenger prostaglandin.[75] Other enzyme inhibitors are poisons. For example, the poison cyanide is an irreversible enzyme inhibitor that combines with the copper and iron in the active site of the enzyme cytochrome c oxidase and blocks cellular respiration.[78]
+
+As enzymes are made up of proteins, their actions are sensitive to change in many physio chemical factors such as pH, temperature, substrate concentration, etc.
+
+The following table shows pH optima for various enzymes.[79]
+
+Enzymes serve a wide variety of functions inside living organisms. They are indispensable for signal transduction and cell regulation, often via kinases and phosphatases.[80] They also generate movement, with myosin hydrolyzing ATP to generate muscle contraction, and also transport cargo around the cell as part of the cytoskeleton.[81] Other ATPases in the cell membrane are ion pumps involved in active transport. Enzymes are also involved in more exotic functions, such as luciferase generating light in fireflies.[82] Viruses can also contain enzymes for infecting cells, such as the HIV integrase and reverse transcriptase, or for viral release from cells, like the influenza virus neuraminidase.[83]
+
+An important function of enzymes is in the digestive systems of animals. Enzymes such as amylases and proteases break down large molecules (starch or proteins, respectively) into smaller ones, so they can be absorbed by the intestines. Starch molecules, for example, are too large to be absorbed from the intestine, but enzymes hydrolyze the starch chains into smaller molecules such as maltose and eventually glucose, which can then be absorbed. Different enzymes digest different food substances. In ruminants, which have herbivorous diets, microorganisms in the gut produce another enzyme, cellulase, to break down the cellulose cell walls of plant fiber.[84]
+
+Several enzymes can work together in a specific order, creating metabolic pathways.[1]:30.1 In a metabolic pathway, one enzyme takes the product of another enzyme as a substrate. After the catalytic reaction, the product is then passed on to another enzyme. Sometimes more than one enzyme can catalyze the same reaction in parallel; this can allow more complex regulation: with, for example, a low constant activity provided by one enzyme but an inducible high activity from a second enzyme.[85]
+
+Enzymes determine what steps occur in these pathways. Without enzymes, metabolism would neither progress through the same steps and could not be regulated to serve the needs of the cell. Most central metabolic pathways are regulated at a few key steps, typically through enzymes whose activity involves the hydrolysis of ATP. Because this reaction releases so much energy, other reactions that are thermodynamically unfavorable can be coupled to ATP hydrolysis, driving the overall series of linked metabolic reactions.[1]:30.1
+
+There are five main ways that enzyme activity is controlled in the cell.[1]:30.1.1
+
+Enzymes can be either activated or inhibited by other molecules. For example, the end product(s) of a metabolic pathway are often inhibitors for one of the first enzymes of the pathway (usually the first irreversible step, called committed step), thus regulating the amount of end product made by the pathways. Such a regulatory mechanism is called a negative feedback mechanism, because the amount of the end product produced is regulated by its own concentration.[86]:141–48 Negative feedback mechanism can effectively adjust the rate of synthesis of intermediate metabolites according to the demands of the cells. This helps with effective allocations of materials and energy economy, and it prevents the excess manufacture of end products. Like other homeostatic devices, the control of enzymatic action helps to maintain a stable internal environment in living organisms.[86]:141
+
+Examples of post-translational modification include phosphorylation, myristoylation and glycosylation.[86]:149–69 For example, in the response to insulin, the phosphorylation of multiple enzymes, including glycogen synthase, helps control the synthesis or degradation of glycogen and allows the cell to respond to changes in blood sugar.[87] Another example of post-translational modification is the cleavage of the polypeptide chain. Chymotrypsin, a digestive protease, is produced in inactive form as chymotrypsinogen in the pancreas and transported in this form to the stomach where it is activated. This stops the enzyme from digesting the pancreas or other tissues before it enters the gut. This type of inactive precursor to an enzyme is known as a zymogen[86]:149–53 or proenzyme.
+
+Enzyme production (transcription and translation of enzyme genes) can be enhanced or diminished by a cell in response to changes in the cell's environment. This form of gene regulation is called enzyme induction. For example, bacteria may become resistant to antibiotics such as penicillin because enzymes called beta-lactamases are induced that hydrolyse the crucial beta-lactam ring within the penicillin molecule.[88] Another example comes from enzymes in the liver called cytochrome P450 oxidases, which are important in drug metabolism. Induction or inhibition of these enzymes can cause drug interactions.[89] Enzyme levels can also be regulated by changing the rate of enzyme degradation.[1]:30.1.1 The opposite of enzyme induction is enzyme repression.
+
+Enzymes can be compartmentalized, with different metabolic pathways occurring in different cellular compartments. For example, fatty acids are synthesized by one set of enzymes in the cytosol, endoplasmic reticulum and Golgi and used by a different set of enzymes as a source of energy in the mitochondrion, through β-oxidation.[90] In addition, trafficking of the enzyme to different compartments may change the degree of protonation (e.g., the neutral cytoplasm and the acidic lysosome) or oxidative state (e.g., oxidizing periplasm or reducing cytoplasm) which in turn affects enzyme activity.[91] In contrast to partitioning into membrane bound organelles, enzyme subcellular localisation may also be altered through polymerisation of enzymes into macromolecular cytoplasmic filaments.[92][93]
+
+In multicellular eukaryotes, cells in different organs and tissues have different patterns of gene expression and therefore have different sets of enzymes (known as isozymes) available for metabolic reactions. This provides a mechanism for regulating the overall metabolism of the organism. For example, hexokinase, the first enzyme in the glycolysis pathway, has a specialized form called glucokinase expressed in the liver and pancreas that has a lower affinity for glucose yet is more sensitive to glucose concentration.[94] This enzyme is involved in sensing blood sugar and regulating insulin production.[95]
+
+Since the tight control of enzyme activity is essential for homeostasis, any malfunction (mutation, overproduction, underproduction or deletion) of a single critical enzyme can lead to a genetic disease. The malfunction of just one type of enzyme out of the thousands of types present in the human body can be fatal. An example of a fatal genetic disease due to enzyme insufficiency is Tay–Sachs disease, in which patients lack the enzyme hexosaminidase.[96][97]
+
+One example of enzyme deficiency is the most common type of phenylketonuria. Many different single amino acid mutations in the enzyme phenylalanine hydroxylase, which catalyzes the first step in the degradation of phenylalanine, result in build-up of phenylalanine and related products. Some mutations are in the active site, directly disrupting binding and catalysis, but many are far from the active site and reduce activity by destabilising the protein structure, or affecting correct oligomerisation.[98][99] This can lead to intellectual disability if the disease is untreated.[100] Another example is pseudocholinesterase deficiency, in which the body's ability to break down choline ester drugs is impaired.[101] 
+Oral administration of enzymes can be used to treat some functional enzyme deficiencies, such as pancreatic insufficiency[102] and lactose intolerance.[103]
+
+Another way enzyme malfunctions can cause disease comes from germline mutations in genes coding for DNA repair enzymes. Defects in these enzymes cause cancer because cells are less able to repair mutations in their genomes. This causes a slow accumulation of mutations and results in the development of cancers. An example of such a hereditary cancer syndrome is xeroderma pigmentosum, which causes the development of skin cancers in response to even minimal exposure to ultraviolet light.[104][105]
+
+Similar to any other protein, enzymes change over time through mutations and sequence divergence. Given their central role in metabolism, enzyme evolution plays a critical role in adaptation. A key question is therefore whether and how enzymes can change their enzymatic activities alongside. It is generally accepted that many new enzyme activities have evolved through gene duplication and mutation of the duplicate copies although evolution can also happen without duplication. One example of an enzyme that has changed its activity is the ancestor of methionyl amino peptidase (MAP) and creatine amidinohydrolase (creatinase) which are clearly homologous but catalyze very different reactions (MAP removes the amino-terminal methionine in new proteins while creatinase hydrolyses creatine to sarcosine and urea). In addition, MAP is metal-ion dependent while creatinase is not, hence this property was also lost over time.[106] Small changes of enzymatic activity are extremely common among enzymes. In particular, substrate binding specificity (see above) can easily and quickly change with single amino acid changes in their substrate binding pockets. This is frequently seen in the main enzyme classes such as kinases.[107]
+
+Artificial (in vitro) evolution is now commonly used to modify enzyme activity or specificity for industrial applications (see below).
+
+Enzymes are used in the chemical industry and other industrial applications when extremely specific catalysts are required. Enzymes in general are limited in the number of reactions they have evolved to catalyze and also by their lack of stability in organic solvents and at high temperatures. As a consequence, protein engineering is an active area of research and involves attempts to create new enzymes with novel properties, either through rational design or in vitro evolution.[108][109] These efforts have begun to be successful, and a few enzymes have now been designed "from scratch" to catalyze reactions that do not occur in nature.[110]
+
+General
+
+Etymology and history
+
+
+
+Enzyme structure and mechanism
+
+Kinetics and inhibition
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+

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+The viola (/viˈoʊlə/ vee-OH-lə,[1][2] also UK: /vaɪˈoʊlə/ vy-OH-lə,[3][4][a] Italian: [ˈvjɔːla, viˈɔːla]) is a string instrument that is bowed, plucked, or played with varying techniques. It is slightly larger than a violin and has a lower and deeper sound. Since the 18th century, it has been the middle or alto voice of the violin family, between the violin (which is tuned a perfect fifth above) and the cello (which is tuned an octave below).[5] The strings from low to high are typically tuned to C3, G3, D4, and A4.
+
+In the past, the viola varied in size and style, as did its names. The word viola originates from the Italian language. The Italians often used the term: viola da braccio meaning literally: 'of the arm'. "Brazzo" was another Italian word for the viola, which the Germans adopted as Bratsche. The French had their own names: cinquiesme was a small viola, haute contre was a large viola, and taile was a tenor. Today, the French use the term alto, a reference to its range.
+
+The viola was popular in the heyday of five-part harmony, up until the eighteenth century, taking three lines of the harmony and occasionally playing the melody line. Music for the viola differs from most other instruments in that it primarily uses the alto clef. When viola music has substantial sections in a higher register, it switches to the treble clef to make it easier to read.
+
+The viola often plays the "inner voices" in string quartets and symphonic writing, and it is more likely than the first violin to play accompaniment parts. The viola occasionally plays a major, soloistic role in orchestral music. Examples include the symphonic poem, "Don Quixote", by Richard Strauss, and the symphony, "Harold en Italie", by Hector Berlioz. In the earlier part of the 20th century, more composers began to write for the viola, encouraged by the emergence of specialized soloists such as Lionel Tertis and William Primrose. English composers Arthur Bliss, York Bowen, Benjamin Dale, Frank Bridge, Benjamin Britten, Rebecca Clarke and Ralph Vaughan Williams all wrote substantial chamber and concert works. Many of these pieces were commissioned by, or written for Lionel Tertis. William Walton, Bohuslav Martinů, Toru Takemitsu, Tibor Serly, Alfred Schnittke, and Béla Bartók have written well-known viola concertos. Paul Hindemith, who was a violist, wrote a substantial amount of music for viola, including the concerto, "Der Schwanendreher". The concerti by Béla Bartók, Paul Hindemith, Carl Stamitz, Georg Philipp Telemann, and William Walton are considered major works of the viola repertoire.
+
+The viola is similar in material and construction to the violin. A full-size viola's body is between 25 mm (1 in) and 100 mm (4 in) longer than the body of a full-size violin (i.e., between 38 and 46 cm [15–18 in]), with an average length of 41 cm (16 in). Small violas typically made for children typically start at 30 cm (12 in), which is equivalent to a half-size violin. For a child who needs a smaller size, a fractional-sized violin is often strung with the strings of a viola.[6] Unlike the violin, the viola does not have a standard full size. The body of a viola would need to measure about 51 cm (20 in) long to match the acoustics of a violin, making it impractical to play in the same manner as the violin.[7] For centuries, viola makers have experimented with the size and shape of the viola, often adjusting proportions or shape to make a lighter instrument with shorter string lengths, but with a large enough sound box to retain the viola sound. Prior to the eighteenth century, violas had no uniform size. Large violas (tenors) were designed to play the lower register viola lines or second viola in five part harmony depending on instrumentation. A smaller viola, nearer the size of the violin, was called an alto viola. It was more suited to higher register writing, as in the viola 1 parts, as their sound was usually richer in the upper register. Its size was not as conducive to a full tone in the lower register.
+
+Several experiments have intended to increase the size of the viola to improve its sound. Hermann Ritter's viola alta, which measured about 48 cm (19 in), was intended for use in Wagner's operas.[8] The Tertis model viola, which has wider bouts and deeper ribs to promote a better tone, is another slightly "nonstandard" shape that allows the player to use a larger instrument. Many experiments with the acoustics of a viola, particularly increasing the size of the body, have resulted in a much deeper tone, making it resemble the tone of a cello. Since many composers wrote for a traditional-sized viola, particularly in orchestral music, changes in the tone of a viola can have unintended consequences upon the balance in ensembles.
+
+One of the most notable makers of violas of the twentieth century was Englishman A. E. Smith, whose violas are sought after and highly valued. Many of his violas remain in Australia, his country of residence, where during some decades the violists of the Sydney Symphony Orchestra had a dozen of them in their section.
+
+More recent (and more radically shaped) innovations have addressed the ergonomic problems associated with playing the viola by making it shorter and lighter, while finding ways to keep the traditional sound. These include the Otto Erdesz "cutaway" viola, which has one shoulder cut out to make shifting easier;[9] the "Oak Leaf" viola, which has two extra bouts; viol-shaped violas such as Joseph Curtin's "Evia" model, which also uses a moveable neck and maple-veneered carbon fibre back, to reduce weight:[10] violas played in the same manner as cellos (see vertical viola); and the eye-catching "Dalí-esque" shapes of both Bernard Sabatier's violas in fractional sizes—which appear to have melted—and David Rivinus' Pellegrina model violas.[11]
+
+Other experiments that deal with the "ergonomics vs. sound" problem have appeared. The American composer Harry Partch fitted a viola with a cello neck to allow the use of his 43-tone scale. Luthiers have also created five-stringed violas, which allow a greater playing range.
+
+A person who plays the viola is called a violist or a viola player. The technique required for playing a viola has certain differences compared with that of a violin, partly because of its larger size: the notes are spread out farther along the fingerboard and often require different fingerings. The viola's less responsive strings and the heavier bow warrant a somewhat different bowing technique, and a violist has to lean more intensely on the strings.[12]
+
+The viola's four strings are normally tuned in fifths: the lowest string is C (an octave below middle C), with G, D and A above it. This tuning is exactly one fifth below the violin,[14] so that they have three strings in common—G, D, and A—and is one octave above the cello.
+
+Each string of a viola is wrapped around a peg near the scroll and is tuned by turning the peg. Tightening the string raises the pitch; loosening the string lowers the pitch. The A string is normally tuned first, typically to a pitch of 440 Hz or 442 Hz. The other strings are then tuned to it in intervals of perfect fifths, sometimes by bowing two strings simultaneously. Most violas also have adjusters—fine tuners, that make finer changes. These adjust the tension of the string via rotating a small knob at the opposite or tailpiece end of the string. Such tuning is generally easier to learn than using the pegs, and adjusters are usually recommended for younger players and put on smaller violas, though pegs and adjusters are usually used together. Adjusters work best, and are most useful, on metal strings. It is common to use one on the A string, which is most prone to breaking, even if the others are not equipped with them. Some violists reverse the stringing of the C and G pegs, so that the thicker C string does not turn so severe an angle over the nut, although this is uncommon.
+
+Small, temporary tuning adjustments can also be made by stretching a string with the hand. A string may be tuned down by pulling it above the fingerboard, or tuned up by pressing the part of the string in the pegbox. These techniques may be useful in performance, reducing the ill effects of an out-of-tune string until an opportunity to tune properly.
+
+The tuning C–G–D–A is used for the great majority of all viola music. However, other tunings are occasionally employed, both in classical music, where the technique is known as scordatura, and in some folk styles. Mozart, in his Sinfonia Concertante for Violin, Viola and Orchestra in E♭, wrote the viola part in D major, and specified that the violist raise the strings in pitch by a semitone. He probably intended to give the viola a brighter tone so the rest of the ensemble wouldn't overpower it. Lionel Tertis, in his transcription of the Elgar cello concerto, wrote the slow movement with the C string tuned down to B♭, enabling the viola to play one passage an octave lower. Occasionally the C string may also be tuned up to D.
+
+A renewal of interest in the viola by performers and composers in the twentieth century led to increased research devoted to the instrument. Paul Hindemith and Vadim Borisovsky made an early attempt at an organization, in 1927, with the Violists' World Union. But it was not until 1968, with the creation of the Viola-Forschungsgesellschaft, now the International Viola Society (IVS), that a lasting organization took hold.[citation needed] The IVS now consists of twelve chapters around the world, the largest being the American Viola Society (AVS), which publishes the Journal of the American Viola Society. In addition to the journal, the AVS sponsors the David Dalton Research Competition and the Primrose International Viola Competition.
+
+The 1960s also saw the beginning of several research publications devoted to the viola, beginning with Franz Zeyringer's, "Literatur für Viola", which has undergone several versions, the most recent being in 1985. In 1980, Maurice Riley produced the first attempt at a comprehensive history of the viola, in his History of the Viola, which was followed with a second volume in 1991. The IVS published the multi-language Viola Yearbook from 1979 to 1994, during which several other national chapters of the IVS published respective newsletters. The Primrose International Viola Archive at Brigham Young University houses the greatest amount of material related to the viola, including scores, recordings, instruments, and archival materials from some of the world's greatest violists.[citation needed]
+
+Music that is written for the viola differs from that of most other instruments, in that it primarily uses the alto clef, which is otherwise rarely used. The trombone occasionally uses the alto clef, but not primarily. (The comparatively rare alto trombone primarily uses the alto clef.) Viola music employs the treble clef when there are substantial sections of music written in a higher register. The alto clef is defined by the placement of C4 on the center line of the staff. In treble clef, this note is placed one line below the staff and in the bass clef (used, notably, by the cello and double bass) it is placed one line above.[15]
+
+As the viola is tuned exactly one octave above the cello (meaning that the viola retains the same string notes as the cello, but an octave up), pieces written for the cello can be easily transposed to the alto clef. For example, there are numerous editions of Bach's Cello Suites transcribed for viola that retain the original key, notes, and musical patterns. The viola also has the advantage of smaller strings, which means that the intervals meant for cello are easier to play on the viola.
+
+In early orchestral music, the viola part was usually limited to filling in harmonies, with very little melodic material assigned to it.  When the viola was given a melodic part, it was often duplicated (or was in unison with) the melody played by other strings.
+
+The concerti grossi, "Brandenburg Concerti", composed by J. S. Bach, were unusual in their use of viola. The third concerto grosso, scored for three violins, three violas, and lower strings with basso continuo, requires occasional virtuosity from the violists. The sixth concerto grosso, "Brandenburg Concerto No. 6", which was scored for 2 violas "concertino", cello, 2 violas da gamba, and continuo, had the two violas playing the primary melodic role.[16] He also used this unusual ensemble in his cantata, Gleichwie der Regen und Schnee vom Himmel fällt, BWV 18 and in Mein Herze schwimmt im Blut, BWV 199, the chorale is accompanied by an obbligato viola.
+
+There are a few Baroque and Classical concerti, such as those by Georg Philipp Telemann (one of the earliest viola concertos known), Alessandro Rolla, Franz Anton Hoffmeister and Carl Stamitz. Hector Berlioz's, "Harold in Italy", was written for solo viola and orchestra.
+
+The viola plays an important role in chamber music. Mozart used the viola in more creative ways when he wrote his six string quintets. The quintets use two violas, which frees them (especially the first viola) for solo passages and increases the variety of writing that is possible for the ensemble. Mozart also wrote for the viola in his, "Sinfonia Concertante", a set of two duets for violin and viola, and the Kegelstatt Trio for viola, clarinet, and piano. The young Felix Mendelssohn wrote a little-known Viola Sonata in C minor (without opus number, but dating from 1824). Robert Schumann wrote his Märchenbilder for viola and piano. He also wrote a set of four pieces for clarinet, viola, and piano, Märchenerzählungen.
+
+Max Bruch wrote a romance for viola and orchestra, his Op. 85, which explores the emotive capabilities of the viola's timbre. In addition, his Eight pieces for clarinet, viola, and piano, Op. 83, features the viola in a very prominent, solo aspect throughout. His Concerto for Clarinet, Viola, and Orchestra, Op. 88 has been quite prominent in the repertoire and has been recorded by prominent violists throughout the 20th century.
+
+From his earliest works, Brahms wrote music that prominently featured the viola. Among his first published pieces of chamber music, the sextets for strings Op. 18 and Op. 36 contain what amounts to solo parts for both violas. Late in life, he wrote two greatly admired sonatas for clarinet and piano, his Op. 120 (1894): he later transcribed these works for the viola (the solo part in his horn trio is also available in a transcription for viola). Brahms also wrote "Two Songs for Alto with Viola and Piano", Op. 91, "Gestillte Sehnsucht" ("Satisfied Longing") and "Geistliches Wiegenlied" ("Spiritual Lullaby") as presents for the famous violinist Joseph Joachim and his wife, Amalie. Dvořák played the viola and apparently said that it was his favorite instrument: his chamber music is rich in important parts for the viola. Another Czech composer, Bedřich Smetana, included a significant viola part in his quartet "From My Life": the quartet begins with an impassioned statement by the viola.  Bach, Mozart and Beethoven all occasionally played the viola part in chamber music.
+
+The viola occasionally has a major role in orchestral music, a prominent example being Richard Strauss' tone poem Don Quixote for solo cello and viola and orchestra. Other examples are the "Ysobel" variation of Edward Elgar's Enigma Variations and the solo in his other work, In the South (Alassio), the pas de deux scene from Act 2 of Adolphe Adam's Giselle and the "La Paix" movement of Léo Delibes's ballet Coppélia, which features a lengthy viola solo.
+
+Gabriel Fauré's Requiem was originally scored (in 1888) with divided viola sections, lacking the usual violin sections, having only a solo violin for the Sanctus. It was later scored for orchestra with violin sections, and published in 1901. Recordings of the older scoring with violas are available.
+
+While the viola repertoire is quite large, the amount written by well-known pre-20th-century composers is relatively small. There are many transcriptions of works for other instruments for the viola and the large number of 20th-century compositions is very diverse. See "The Viola Project" at the San Francisco Conservatory of Music, where Professor of Viola Jodi Levitz has paired a composer with each of her students, resulting in a recital of brand-new works played for the very first time.
+
+In the earlier part of the 20th century, more composers began to write for the viola, encouraged by the emergence of specialized soloists such as Lionel Tertis. Englishmen Arthur Bliss, York Bowen, Benjamin Dale, and Ralph Vaughan Williams all wrote chamber and concert works for Tertis. William Walton, Bohuslav Martinů, and Béla Bartók wrote well-known viola concertos. Paul Hindemith wrote a substantial amount of music for the viola; being himself a violist, he often performed his own works. Claude Debussy's Sonata for flute, viola and harp has inspired a significant number of other composers to write for this combination.
+
+Charles Wuorinen composed his virtuosic Viola Variations in 2008 for Lois Martin. Elliott Carter also wrote several works for viola including his Elegy (1943) for viola and piano; it was subsequently transcribed for clarinet. Ernest Bloch, a Swiss-born American composer best known for his compositions inspired by Jewish music, wrote two famous works for viola, the Suite 1919 and the Suite Hébraïque for solo viola and orchestra. Rebecca Clarke was a 20th-century composer and violist who also wrote extensively for the viola. Lionel Tertis records that Edward Elgar (whose cello concerto Tertis transcribed for viola, with the slow movement in scordatura), Alexander Glazunov (who wrote an Elegy, Op. 44, for viola and piano), and Maurice Ravel all promised concertos for viola, yet all three died before doing any substantial work on them.
+
+In the latter part of the 20th century a substantial repertoire was produced for the viola; many composers including Miklós Rózsa, Revol Bunin, Alfred Schnittke, Sofia Gubaidulina, Giya Kancheli and Krzysztof Penderecki, have written viola concertos. The American composer Morton Feldman wrote a series of works entitled The Viola in My Life, which feature concertante viola parts. In spectral music, the viola has been sought after because of its lower overtone partials that are more easily heard than on the violin. Spectral composers like Gérard Grisey, Tristan Murail, and Horațiu Rădulescu have written solo works for viola. Neo-Romantic, post-Modern composers have also written significant works for viola including Robin Holloway Viola Concerto op.56 and Sonata op.87, and Peter Seabourne a large five movement work with piano, Pietà.
+
+The viola is sometimes used in contemporary popular music, mostly in the avant-garde. John Cale of The Velvet Underground used the viola, as do some modern groups such as alternative rock band 10,000 Maniacs, Imagine Dragons, folk duo John & Mary, Flobots, British Sea Power, Quargs (Mya) Greene of Love Ghost and others. Jazz music has also seen its share of violists, from those used in string sections in the early 1900s to a handful of quartets and soloists emerging from the 1960s onward. It is quite unusual though, to use individual bowed string instruments in contemporary popular music.
+
+Although not as commonly used as the violin in folk music, the viola is nevertheless used by many folk musicians across the world. Extensive research into the historical and current use of the viola in folk music has been carried out by Dr. Lindsay Aitkenhead. Players in this genre include Eliza Carthy, Mary Ramsey, Helen Bell, and Nancy Kerr. Clarence "Gatemouth" Brown was the viola's most prominent exponent in the genre of blues.
+
+The viola is also an important accompaniment instrument in Slovakian, Hungarian and Romanian folk string band music, especially in Transylvania. Here the instrument has three strings tuned G3–D4–A3 (note that the A is an octave lower than found on the standard instrument), and the bridge is flattened with the instrument playing chords in a strongly rhythmic manner. In this usage, it is called a kontra or brácsa (pronounced "bra-cha", from German Bratsche, "viola").
+
+There are few well-known viola virtuoso soloists, perhaps because little virtuoso viola music was written before the twentieth century. Pre-twentieth century viola players of note include Carl Stamitz, Alessandro Rolla, Antonio Rolla, Chrétien Urhan, Casimir Ney, Louis van Waefelghem, and Hermann Ritter. Important viola pioneers from the twentieth century were Lionel Tertis, William Primrose, composer/performer Paul Hindemith, Théophile Laforge, Cecil Aronowitz, Maurice Vieux, Vadim Borisovsky, Lillian Fuchs, Dino Asciolla, Frederick Riddle, Walter Trampler, Ernst Wallfisch, Csaba Erdélyi, the only violist to ever win the Carl Flesch International Violin Competition, and Emanuel Vardi, the first violist to record the 24 Caprices by Paganini on viola. Many noted violinists have publicly performed and recorded on the viola as well, among them Eugène Ysaÿe, Yehudi Menuhin, David Oistrakh, Pinchas Zukerman, Maxim Vengerov, Julian Rachlin and Nigel Kennedy.
+
+Among the great composers, several preferred the viola to the violin when they were playing in ensembles,[18] the most noted being Ludwig van Beethoven, Johann Sebastian Bach[19] and Wolfgang Amadeus Mozart. Other composers also chose to play the viola in ensembles, including Joseph Haydn, Franz Schubert, Felix Mendelssohn, Antonín Dvořák, and Benjamin Britten. Among those noted both as violists and as composers are Rebecca Clarke and Paul Hindemith. Contemporary composers and violists Kenji Bunch, Scott Slapin, and Lev Zhurbin have written a number of works for viola.
+
+Amplification of a viola with a pickup, an instrument amplifier (and speaker), and adjusting the tone with a graphic equalizer can make up for the comparatively weaker output of a violin-family instrument string tuned to notes below G3. There are two types of instruments used for electric viola: regular acoustic violas fitted with a piezoelectric pickup and specialized electric violas, which have little or no body. While traditional acoustic violas are typically only available in historically used earth tones (e.g., brown, reddish brown, blonde), electric violas may be traditional colors or they may use bright colors, such as red, blue or green. Some electric violas are made of materials other than wood.
+
+Most electric instruments with lower strings are violin-sized, as they use the amp and speaker to create a big sound, so they don't need a large soundbox. Indeed, some electric violas have little or no soundbox, and thus rely on entirely on amplification. Fewer electric violas are available relative to electric violins. it can be hard for violists who prefer physical size or familiar touch references of a viola-sized instrument, when they must use an electric viola that uses a smaller violin-sized body. Welsh musician John Cale, formerly of The Velvet Underground, is one of the more famous users of such an electric viola, who has used them both for melodies in his solo work and for drones in his work with The Velvet Underground (e.g. "Venus in Furs"). Other notable players of the electric viola are Geoffrey Richardson of Caravan and Mary Ramsey of 10,000 Maniacs.
+
+Instruments may be built with an internal preamplifier, or may put out an unbuffered transducer signal. While such signals may be fed directly to an amplifier or mixing board, they often benefit from an external preamp/equalizer on the end of a short cable, before being fed to the sound system. In rock and other loud styles, the electric viola player may use effects units such as reverb or overdrive.

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+
+
+A wind turbine, or alternatively referred to as a wind energy converter, is a device that converts the wind's kinetic energy into electrical energy.
+
+Wind turbines are manufactured  in a wide range of vertical and horizontal axis. The smallest turbines are used for applications such as battery charging for auxiliary power for boats or caravans or to power traffic warning signs. Larger turbines can be used for making contributions to a domestic power supply while selling unused power back to the utility supplier via the electrical grid. Arrays of large turbines, known as wind farms, are becoming an increasingly important source of intermittent  renewable energy and are used by many countries as part of a strategy to reduce their reliance on fossil fuels. One assessment claimed that, as of 2009[update], wind had the "lowest relative greenhouse gas emissions, the least water consumption demands and... the most favourable social impacts" compared to photovoltaic, hydro, geothermal, coal and gas.[1]
+
+The windwheel of Hero of Alexandria (10 AD – 70 AD) marks one of the first recorded instances of wind powering a machine in history.[2][3] However, the first known practical wind power plants were built in Sistan, an Eastern province of Persia (now Iran), from the 7th century. These "Panemone" were vertical axle windmills, which had long vertical drive shafts with rectangular blades.[4] Made of six to twelve sails covered in reed matting or cloth material, these windmills were used to grind grain or draw up water, and were used in the gristmilling and sugarcane industries.[5]
+
+Wind power first appeared in Europe during the Middle Ages. The first historical records of their use in England date to the 11th or 12th centuries, there are reports of German crusaders taking their windmill-making skills to Syria around 1190.[6] By the 14th century, Dutch windmills were in use to drain areas of the Rhine delta. Advanced wind turbines were described by Croatian inventor Fausto Veranzio. In his book Machinae Novae (1595) he described vertical axis wind turbines with curved or V-shaped blades.
+
+The first electricity-generating wind turbine was a battery charging machine installed in July 1887 by Scottish academic James Blyth to light his holiday home in Marykirk, Scotland.[7] Some months later American inventor Charles F. Brush was able to build the first automatically operated wind turbine after consulting local University professors and colleagues Jacob S. Gibbs and Brinsley Coleberd and successfully getting the blueprints peer-reviewed for electricity production in Cleveland, Ohio.[7] Although Blyth's turbine was considered uneconomical in the United Kingdom,[7] electricity generation by wind turbines was more cost effective in countries with widely scattered populations.[6]
+
+In Denmark by 1900, there were about 2500 windmills for mechanical loads such as pumps and mills, producing an estimated combined peak power of about 30 MW. The largest machines were on 24-meter (79 ft) towers with four-bladed 23-meter (75 ft) diameter rotors. By 1908, there were 72 wind-driven electric generators operating in the United States from 5 kW to 25 kW. Around the time of World War I, American windmill makers were producing 100,000 farm windmills each year, mostly for water-pumping.[9]
+
+By the 1930s, wind generators for electricity were common on farms, mostly in the United States where distribution systems had not yet been installed. In this period, high-tensile steel was cheap, and the generators were placed atop prefabricated open steel lattice towers.
+
+A forerunner of modern horizontal-axis wind generators was in service at Yalta, USSR in 1931. This was a 100 kW generator on a 30-meter (98 ft) tower, connected to the local 6.3 kV distribution system. It was reported to have an annual capacity factor of 32 percent, not much different from current wind machines.[10][11]
+
+In the autumn of 1941, the first megawatt-class wind turbine was synchronized to a utility grid in Vermont. The Smith–Putnam wind turbine only ran for 1,100 hours before suffering a critical failure. The unit was not repaired, because of a shortage of materials during the war.
+
+The first utility grid-connected wind turbine to operate in the UK was built by John Brown & Company in 1951 in the Orkney Islands.[7][12]
+
+Despite these diverse developments, developments in fossil fuel systems almost entirely eliminated any wind turbine systems larger than supermicro size.  In the early 1970s, however, anti-nuclear protests in Denmark spurred artisan mechanics to develop microturbines of 22 kW.  Organizing owners into associations and co-operatives lead to the lobbying of the government and utilities and provided incentives for larger turbines throughout the 1980s and later.  Local activists in Germany, nascent turbine manufacturers in Spain, and large investors in the United States in the early 1990s then lobbied for policies that stimulated the industry in those countries.
+
+It has been argued that expanding use of wind power will lead to increasing geopolitical competition over critical materials for wind turbines such as rare earth elements neodymium, praseodymium, and dysprosium. But this perspective has been criticised for failing to recognise that most wind turbines do not use permanent magnets and for underestimating the power of economic incentives for expanded production of these minerals.[13]
+
+Wind Power Density (WPD) is a quantitative measure of  wind energy available at any location. It is the mean annual power available per square meter of swept area of a turbine, and is calculated for different heights above ground. Calculation of wind power density includes the effect of wind velocity and air density.[14]
+
+Wind turbines are classified by the wind speed they are designed for, from class I to class III, with A to C referring to the turbulence intensity of the wind.[15]
+
+Conservation of mass requires that the amount of air entering and exiting a turbine must be equal. Accordingly, Betz's law gives the maximal achievable extraction of wind power by a wind turbine as 16/27 (59.3%) of the rate at which the kinetic energy of the air arrives at the turbine.[16]
+
+The maximum theoretical power output of a wind machine is thus 16/27 times the rate at which kinetic energy of the air arrives at the effective disk area of the machine. If the effective area of the disk is A, and the wind velocity v, the maximum theoretical power output P is:
+
+where ρ is the air density.
+
+Wind-to-rotor efficiency (including rotor blade friction and drag) are among the factors affecting the final price of wind power.[17]
+Further inefficiencies, such as gearbox losses, generator and converter losses, reduce the power delivered by a wind turbine. To protect components from undue wear, extracted power is held constant above the rated operating speed as theoretical power increases at the cube of wind speed, further reducing theoretical efficiency. In 2001, commercial utility-connected turbines delivered 75% to 80% of the Betz limit of power extractable from the wind, at rated operating speed.[18][19][needs update]
+
+Efficiency can decrease slightly over time, one of the main reasons being dust and insect carcasses on the blades which alters the aerodynamic profile and essentially reduces the lift to drag ratio of the airfoil. Analysis of 3128 wind turbines older than 10 years in Denmark showed that half of the turbines had no decrease, while the other half saw a production decrease of 1.2% per year.[20]  Ice accretion on turbine blades has also been found to greatly reduce the efficiency of wind turbines, which is a common challenge in cold climates where in-cloud icing and freezing rain events occur.[21] Vertical turbine designs have much lower efficiency than standard horizontal designs.[22]
+
+In general, more stable and constant weather conditions (most notably wind speed) result in an average of 15% greater efficiency than that of a wind turbine in unstable weather conditions, thus allowing up to a 7% increase in wind speed under stable conditions. This is due to a faster recovery wake and greater flow entrainment that occur in conditions of higher atmospheric stability. However, wind turbine wakes have been found to recover faster under unstable atmospheric conditions as opposed to a stable environment.[23]
+
+Different materials have been found to have varying effects on the efficiency of wind turbines. In an Ege University experiment, three wind turbines (Each with three blades with diameters of one meter) were constructed with blades made of different materials: A glass and glass/carbon epoxy, glass/carbon, and glass/polyester. When tested, the results showed that the materials with higher overall masses had a greater friction moment and thus a lower power coefficient.[24]
+
+Wind turbines can rotate about either a horizontal or a vertical axis, the former being both older and more common.[25] They can also include blades, or be bladeless.[26] Vertical designs produce less power and are less common.[27]
+
+Large three-bladed horizontal-axis wind turbines (HAWT) with the blades upwind of the tower produce the overwhelming majority of wind power in the world today. These turbines have the main rotor shaft and electrical generator at the top of a tower, and must be pointed into the wind. Small turbines are pointed by a simple wind vane, while large turbines generally use a wind sensor coupled with a yaw system. Most have a gearbox, which turns the slow rotation of the blades into a quicker rotation that is more suitable to drive an electrical generator.[28] Some turbines use a different type of generator suited to slower rotational speed input. These don't need a gearbox and are called direct-drive, meaning they couple the rotor directly to the generator with no gearbox in between. While permanent magnet direct-drive generators can be more costly due to the rare earth materials required, these gearless turbines are sometimes preferred over gearbox generators because they "eliminate the gear-speed increaser, which is susceptible to significant accumulated fatigue torque loading, related reliability issues, and maintenance costs."[29] There is also the pseudo direct drive mechanism, which has some advantages over the permanent magnet direct drive mechanism.[30][31]
+
+Most horizontal axis turbines have their rotors upwind of the supporting tower. Downwind machines have been built, because they don't need an additional mechanism for keeping them in line with the wind. In high winds, the blades can also be allowed to bend, which reduces their swept area and thus their wind resistance. Despite these advantages, upwind designs are preferred, because the change in loading from the wind as each blade passes behind the supporting tower can cause damage to the turbine.
+
+Turbines used in wind farms for commercial production of electric power are usually three-bladed. These have low torque ripple, which contributes to good reliability. The blades are usually colored white for daytime visibility by aircraft and range in length from 20 to 80 meters (66 to 262 ft). The size and height of turbines increase year by year. Offshore wind turbines are built up to 8 MW today and have a blade length up to 80 meters (260 ft). Designs with 10 to 12 MW are in preparation.[32] Usual multi megawatt turbines have tubular steel towers with a height of 70 m to 120 m and in extremes up to 160 m.
+
+Vertical-axis wind turbines (or VAWTs) have the main rotor shaft arranged vertically. One advantage of this arrangement is that the turbine does not need to be pointed into the wind to be effective, which is an advantage on a site where the wind direction is highly variable. It is also an advantage when the turbine is integrated into a building because it is inherently less steerable.  Also, the generator and gearbox can be placed near the ground, using a direct drive from the rotor assembly to the ground-based gearbox, improving accessibility for maintenance. However, these designs produce much less energy averaged over time, which is a major drawback.[27][33]
+
+The key disadvantages include the relatively low rotational speed with the consequential higher torque and hence higher cost of the drive train, the inherently lower power coefficient, the 360-degree rotation of the aerofoil within the wind flow during each cycle and hence the highly dynamic loading on the blade, the pulsating torque generated by some rotor designs on the drive train, and the difficulty of modelling the wind flow accurately and hence the challenges of analysing and designing the rotor prior to fabricating a prototype.[34]
+
+When a turbine is mounted on a rooftop the building generally redirects wind over the roof and this can double the wind speed at the turbine. If the height of a rooftop mounted turbine tower is approximately 50% of the building height it is near the optimum for maximum wind energy and minimum wind turbulence. While wind speeds within the built environment are generally much lower than at exposed rural sites,[35][36] noise may be a concern and an existing structure may not adequately resist the additional stress.
+
+Subtypes of the vertical axis design include:
+
+"Eggbeater" turbines, or Darrieus turbines, were named after the French inventor, Georges Darrieus.[37] They have good efficiency, but produce large torque ripple and cyclical stress on the tower, which contributes to poor reliability. They also generally require some external power source, or an additional Savonius rotor to start turning, because the starting torque is very low. The torque ripple is reduced by using three or more blades, which results in greater solidity of the rotor. Solidity is measured by blade area divided by the rotor area. Newer Darrieus type turbines are not held up by guy-wires but have an external superstructure connected to the top bearing.[38]
+
+A subtype of Darrieus turbine with straight, as opposed to curved, blades. The cycloturbine variety has variable pitch to reduce the torque pulsation and is self-starting.[39] The advantages of variable pitch are: high starting torque; a wide, relatively flat torque curve; a higher coefficient of performance; more efficient operation in turbulent winds; and a lower blade speed ratio which lowers blade bending stresses. Straight, V, or curved blades may be used.[40]
+
+These are drag-type devices with two (or more) scoops that are used in anemometers, Flettner vents (commonly seen on bus and van roofs), and in some high-reliability low-efficiency power turbines. They are always self-starting if there are at least three scoops.
+
+Twisted Savonius is a modified savonius, with long helical scoops to provide smooth torque. This is often used as a rooftop wind turbine and has even been adapted for ships.[41]
+
+The parallel turbine is similar to the crossflow fan or centrifugal fan. It uses the ground effect. Vertical axis turbines of this type have been tried for many years: a unit producing 10 kW was built by Israeli wind pioneer Bruce Brill in the 1980s.[42][unreliable source?]
+
+Wind turbine design is a careful balance of cost, energy output, and fatigue life.
+
+Wind turbines convert wind energy to electrical energy for distribution. Conventional horizontal axis turbines can be divided into three components:
+
+A 1.5 (MW) wind turbine of a type frequently seen in the United States has a tower 80 meters (260 ft) high. The rotor assembly (blades and hub) weighs 22,000 kilograms (48,000 lb). The nacelle, which contains the generator, weighs 52,000 kilograms (115,000 lb). The concrete base for the tower is constructed using 26,000 kilograms (58,000 lb) reinforcing steel and contains 190 cubic meters (250 cu yd) of concrete. The base is 15 meters (50 ft) in diameter and 2.4 meters (8 ft) thick near the center.[48]
+
+Due to data transmission problems, structural health monitoring of wind turbines is usually performed using several accelerometers and strain gages attached to the nacelle to monitor the gearbox and equipment. Currently, digital image correlation and stereophotogrammetry are used to measure dynamics of wind turbine blades. These methods usually measure displacement and strain to identify location of defects. Dynamic characteristics of non-rotating wind turbines have been measured using digital image correlation and photogrammetry.[49] Three dimensional point tracking has also been used to measure rotating dynamics of wind turbines.[50]
+
+Wind turbine rotor blades are being made longer to increase efficiency. This requires them to be stiff, strong, light and resistant to fatigue.[51] Materials with these properties are composites such as polyester and epoxy, while glass fiber and carbon fiber have been used for the reinforcing.[52] Construction may use manual layup or injection molding.
+
+Companies seek ways to draw greater efficiency from their designs. A predominant way has been to increase blade length and thus rotor diameter. Retrofitting existing turbines with larger blades reduces the work and risks of redesigning the system. The current longest blade is 88.4 m (from LM Wind Power), but by 2021 offshore turbines are expected to be 10-MW with 100 m blades. Longer blades need to be stiffer to avoid deflection, which requires materials with higher stiffness-to-weight ratio. Because the blades need to function over a 100 million load cycles over a period of 20–25 years, the fatigue of the blade materials is also critical.
+
+Materials commonly used in wind turbine blades are described below.
+
+The stiffness of composites is determined by the stiffness of fibers and their volume content. Typically, E-glass fibers are used as main reinforcement in the composites. Typically, the glass/epoxy composites for wind turbine blades contain up to 75% glass by weight. This increases the stiffness, tensile and compression strength. A promising composite material is glass fiber with modified compositions like S-glass, R-glass etc. Other glass fibers developed by Owens Corning are ECRGLAS, Advantex and WindStrand.[53]
+
+Carbon fiber has more tensile strength, higher stiffness and lower density than glass fiber. An ideal candidate for these properties is the spar cap, a structural element of a blade which experiences high tensile loading.[52] A 100-m glass fiber blade could weigh up to 50 metric tons, while using carbon fiber in the spar saves 20% to 30% weight, about 15 metric tons.[54] However, because carbon fiber is ten times more expensive, glass fiber is still dominant.
+
+Instead of making wind turbine blade reinforcements from pure glass or pure carbon, hybrid designs trade weight for cost. For example, for an 8 m blade, a full replacement by carbon fiber would save 80% of weight but increase costs by 150%, while a 30% replacement would save 50% of weight and increase costs by 90%. Hybrid reinforcement materials include E-glass/carbon, E-glass/aramid. The current longest blade by LM Wind Power is made of carbon/glass hybrid composites. More research is needed about the optimal composition of materials [55]
+
+Additions of small amount (0.5 weight %) of nanoreinforcement (carbon nanotubes or nanoclay) in the polymer matrix of composites, fiber sizing or interlaminar layers can improve fatigue resistance, shear or compressive strength, and fracture toughness of the composites by 30% to 80%. Research has also shown that incorporating small amounts of carbon nanotubes (CNT) can increase the lifetime up to 1500%.
+
+As of 2019[update], a wind turbine may cost around $1 million per megawatt.[56]
+
+For the wind turbine blades, while the material cost is much higher for hybrid glass/carbon fiber blades than all-glass fiber blades, labor costs can be lower.  Using carbon fiber allows simpler designs that use less raw material.  The chief manufacturing process in blade fabrication is the layering of plies.  Thinner blades allow reducing the number of layers and so the labor, and in some cases, equate to the cost of labor for glass fiber blades.[57]
+
+Wind turbine parts other than the rotor blades (including the rotor hub, gearbox, frame, and tower) are largely made of steel. Smaller turbines (as well as megawatt-scale Enercon turbines) have begun using aluminum alloys for these components to make turbines lighter and more efficient. This trend may grow if fatigue and strength properties can be improved. 
+Pre-stressed concrete has been increasingly used for the material of the tower, but still requires much reinforcing steel to meet the strength requirement of the turbine. Additionally, step-up gearboxes are being increasingly replaced with variable speed generators, which requires magnetic materials.[51] In particular, this would require an greater supply of the rare earth metal neodymium.
+
+Modern turbines use a couple of tons of copper for generators, cables and such.[58] As of 2018[update], global production of wind turbines use 450,000 tonnes of copper per year.[59]
+
+A study of the material consumption trends and requirements for wind energy in Europe found that bigger turbines have a higher consumption of precious metals but lower material input per kW generated. The current material consumption and stock was compared to input materials for various onshore system sizes. In all EU countries the estimates for 2020 doubled the values consumed in 2009. These countries would need to expand their resources to meet the estimated demand for 2020. For example, currently the EU has 3% of world supply of fluorspar and it requires 14% by 2020. Globally, the main exporting countries are South Africa, Mexico and China.  This is similar with other critical and valuable materials required for energy systems such as magnesium, silver and indium. The levels of recycling of these materials are very low and focusing on that could alleviate supply. Because most of these valuable materials are also used  in other emerging technologies, like light emitting diodes (LEDs), photo voltaics (PVs) and liquid crystal displays (LCDs), their demand is expected to grow.[60]
+
+A study by the United States Geological Survey estimated resources required to fulfill the US commitment to supplying 20% of its electricity from wind power by 2030. It did not consider requirements for small turbines or offshore turbines because those were not common in 2008 when the study was done. Common materials such as cast iron, steel and concrete would increase by 2%–3% compared to 2008. Between 110,000 and 115,000 metric tons of fiber glass would be required per year, a 14% increase. Rare metal use would not increase much compared to available supply, however rare metals that are also used for other technologies such as batteries which are increasing its global demand need to be taken into account. Land required would be 50,000 square kilometers onshore and 11,000 offshore. This would not be a problem in the US due to its vast area and because the same land can be used for farming. A greater challenge would be the variability and transmission to areas of high demand.[61]
+
+Permanent magnets for wind turbine generators contain rare metals such as neodymium (Nd), praseodymium (Pr), Terbium (Tb) and dysprosium (Dy). Systems that use magnetic direct drive turbines require greater amounts of rare metals. Therefore, an increase in wind turbine manufacture would increase the demand for these resources. By 2035, the demand for Nd is estimated to increase by 4,000 to 18,000 tons and for Dy by 200 to 1200 tons. These values are a quarter to half of current production. However, these estimates are very uncertain because technologies are developing rapidly.[62]
+
+Reliance on rare earth minerals for components has risked expense and price volatility as China has been main producer of rare earth minerals (96% in 2009) and was reducing its export quotas.[63] However, in recent years other producers have increased production and China has increased export quotas, leading to a higher supply and lower cost, and a greater viability of large scale use of variable-speed generators.[64]
+
+Glass fiber is the most common material for reinforcement.  Its demand has grown due to growth in construction, transportation and wind turbines. Its global market might reach US$17.4 billion by 2024, compared to US$8.5 billion in 2014. In 2014, Asia Pacific produced more than 45% of the market; now China is the largest producer. The industry receives subsidies from the Chinese government allowing it to export cheaper to the US and Europe. However, price wars have led to anti-dumping measures such as tariffs on Chinese glass fiber.[65]
+
+Interest in recycling blades varies in different markets and depends on the waste legislation and local economics. A challenge in recycling blades is related to the composite material, which is made of a thermosetting matrix and glass fibers or a combination of glass and carbon fibers. Thermosetting matrix cannot be remolded to form new composites. So the options are either to send the blade to landfill, to reuse the blade and the composite material elements found in the blade, or to transform the composite material into a new source of material. In Germany, wind turbine blades are commercially recycled as part of an alternative fuel mix for a cement factory. In the USA the town of Casper, Wyoming has buried 1,000 non-recyclable blades in its landfill site, earning $675,000 for the town. It pointed out that wind farm waste is less toxic than other garbage. Wind turbine blades represent a “vanishingly small fraction” of overall waste in the US, according to the American Wind Energy Association.[66]
+
+A few localities have exploited the attention-getting nature of wind turbines by placing them on public display, either with visitor centers around their bases, or with viewing areas farther away.[67] The wind turbines are generally of conventional horizontal-axis, three-bladed design, and generate power to feed electrical grids, but they also serve the unconventional roles of technology demonstration, public relations, and education.
+
+Small wind turbines may be used for a variety of applications including on- or off-grid residences, telecom towers, offshore platforms, rural schools and clinics, remote monitoring and other purposes that require energy where there is no electric grid, or where the grid is unstable. Small wind turbines may be as small as a fifty-watt generator for boat or caravan use. Hybrid solar and wind powered units are increasingly being used for traffic signage, particularly in rural locations, as they avoid the need to lay long cables from the nearest mains connection point.[68] The U.S. Department of Energy's National Renewable Energy Laboratory (NREL) defines small wind turbines as those smaller than or equal to 100 kilowatts.[69] Small units often have direct drive generators, direct current output, aeroelastic blades, lifetime bearings and use a vane to point into the wind.
+
+Larger, more costly turbines generally have geared power trains, alternating current output, and flaps, and are actively pointed into the wind. Direct drive generators and aeroelastic blades for large wind turbines are being researched.
+
+On most horizontal wind turbine farms, a spacing of about 6–10 times the rotor diameter is often upheld. However, for large wind farms distances of about 15 rotor diameters should be more economical, taking into account typical wind turbine and land costs. This conclusion has been reached by research[70] conducted by Charles Meneveau of Johns Hopkins University[71] and Johan Meyers of Leuven University in Belgium, based on computer simulations[72] that take into account the detailed interactions among wind turbines (wakes) as well as with the entire turbulent atmospheric boundary layer.
+
+Recent research by John Dabiri of Caltech suggests that vertical wind turbines may be placed much more closely together so long as an alternating pattern of rotation is created allowing blades of neighbouring turbines to move in the same direction as they approach one another.[73]
+
+Wind turbines need regular maintenance to stay reliable and available. In the best case turbines are available to generate energy 98% of the time.[74][75]
+
+Modern turbines usually have a small onboard crane for hoisting maintenance tools and minor components. However, large, heavy components like generator, gearbox, blades, and so on are rarely replaced, and a heavy lift external crane is needed in those cases. If the turbine has a difficult access road, a containerized crane can be lifted up by the internal crane to provide heavier lifting.[76]
+
+Installation of new wind turbines can be controversial. An alternative is repowering, where existing wind turbines are replaced with bigger, more powerful ones, sometimes in smaller numbers while keeping or increasing capacity.
+
+Older turbines were in some early cases not required to be removed when reaching the end of their life. Some still stand, waiting to be recycled or repowered.[77][78]
+
+A demolition industry develops to recycle offshore turbines at a cost of DKK 2–4 million per (MW), to be guaranteed by the owner.[79]
+
+Wind turbines produce electricity at between two and six cents per kilowatt hour, which is one of the lowest-priced renewable energy sources.[80][81] As technology needed for wind turbines continued to improve, the prices decreased as well. In addition, there is currently no competitive market for wind energy, because wind is a freely available natural resource, most of which is untapped.[80] The main cost of small wind turbines is the purchase and installation process, which averages between $48,000 and $65,000 per installation. The energy harvested from the turbine will offset the installation cost, as well as provide virtually free energy for years.[82]
+
+Wind turbines provide a clean energy source, use little water,[1] emitting no greenhouse gases and no waste products. Over 1,500 tons of carbon dioxide per year can be eliminated by using a one-megawatt turbine instead of one megawatt of energy from a fossil fuel.[83]
+
+Wind turbines can be very large, reaching over 140 m (460 ft) tall and with blades 55 m (180 ft) long,[84] and people have often complained about their visual impact.
+
+Environmental impact of wind power includes effect on wildlife, but can be mitigated if proper monitoring and mitigation strategies are implemented.[85]  Thousands of birds, including rare species, have been killed by the blades of wind turbines,[86] though wind turbines contribute relatively insignificantly to anthropogenic avian mortality. Wind farms and nuclear power stations are responsible for between 0.3 and 0.4 bird deaths per gigawatt-hour (GWh) of electricity while fossil fueled power stations are responsible for about 5.2 fatalities per GWh. In 2009, for every bird killed by a wind turbine in the US, nearly 500,000 were killed by cats and another 500,000 by buildings.[87] In comparison, conventional coal fired generators contribute significantly more to bird mortality, by incineration when caught in updrafts of smoke stacks and by poisoning with emissions byproducts (including particulates and heavy metals downwind of flue gases). Further, marine life is affected by water intakes of steam turbine cooling towers (heat exchangers) for nuclear and fossil fuel generators, by coal dust deposits in marine ecosystems (e.g. damaging Australia's Great Barrier Reef) and by water acidification from combustion monoxides.
+
+Energy harnessed by wind turbines is intermittent, and is not a "dispatchable" source of power; its availability is based on whether the wind is blowing, not whether electricity is needed. Turbines can be placed on ridges or bluffs to maximize the access of wind they have, but this also limits the locations where they can be placed.[80] In this way, wind energy is not a particularly reliable source of energy. However, it can form part of the energy mix, which also includes power from other sources. Notably, the relative available output from wind and solar sources is often inversely proportional (balancing)[citation needed]. Technology is also being developed to store excess energy, which can then make up for any deficits in supplies.
+
+See also List of most powerful wind turbines

en/1781.html.txt

@@ -0,0 +1,145 @@
+
+
+A wind turbine, or alternatively referred to as a wind energy converter, is a device that converts the wind's kinetic energy into electrical energy.
+
+Wind turbines are manufactured  in a wide range of vertical and horizontal axis. The smallest turbines are used for applications such as battery charging for auxiliary power for boats or caravans or to power traffic warning signs. Larger turbines can be used for making contributions to a domestic power supply while selling unused power back to the utility supplier via the electrical grid. Arrays of large turbines, known as wind farms, are becoming an increasingly important source of intermittent  renewable energy and are used by many countries as part of a strategy to reduce their reliance on fossil fuels. One assessment claimed that, as of 2009[update], wind had the "lowest relative greenhouse gas emissions, the least water consumption demands and... the most favourable social impacts" compared to photovoltaic, hydro, geothermal, coal and gas.[1]
+
+The windwheel of Hero of Alexandria (10 AD – 70 AD) marks one of the first recorded instances of wind powering a machine in history.[2][3] However, the first known practical wind power plants were built in Sistan, an Eastern province of Persia (now Iran), from the 7th century. These "Panemone" were vertical axle windmills, which had long vertical drive shafts with rectangular blades.[4] Made of six to twelve sails covered in reed matting or cloth material, these windmills were used to grind grain or draw up water, and were used in the gristmilling and sugarcane industries.[5]
+
+Wind power first appeared in Europe during the Middle Ages. The first historical records of their use in England date to the 11th or 12th centuries, there are reports of German crusaders taking their windmill-making skills to Syria around 1190.[6] By the 14th century, Dutch windmills were in use to drain areas of the Rhine delta. Advanced wind turbines were described by Croatian inventor Fausto Veranzio. In his book Machinae Novae (1595) he described vertical axis wind turbines with curved or V-shaped blades.
+
+The first electricity-generating wind turbine was a battery charging machine installed in July 1887 by Scottish academic James Blyth to light his holiday home in Marykirk, Scotland.[7] Some months later American inventor Charles F. Brush was able to build the first automatically operated wind turbine after consulting local University professors and colleagues Jacob S. Gibbs and Brinsley Coleberd and successfully getting the blueprints peer-reviewed for electricity production in Cleveland, Ohio.[7] Although Blyth's turbine was considered uneconomical in the United Kingdom,[7] electricity generation by wind turbines was more cost effective in countries with widely scattered populations.[6]
+
+In Denmark by 1900, there were about 2500 windmills for mechanical loads such as pumps and mills, producing an estimated combined peak power of about 30 MW. The largest machines were on 24-meter (79 ft) towers with four-bladed 23-meter (75 ft) diameter rotors. By 1908, there were 72 wind-driven electric generators operating in the United States from 5 kW to 25 kW. Around the time of World War I, American windmill makers were producing 100,000 farm windmills each year, mostly for water-pumping.[9]
+
+By the 1930s, wind generators for electricity were common on farms, mostly in the United States where distribution systems had not yet been installed. In this period, high-tensile steel was cheap, and the generators were placed atop prefabricated open steel lattice towers.
+
+A forerunner of modern horizontal-axis wind generators was in service at Yalta, USSR in 1931. This was a 100 kW generator on a 30-meter (98 ft) tower, connected to the local 6.3 kV distribution system. It was reported to have an annual capacity factor of 32 percent, not much different from current wind machines.[10][11]
+
+In the autumn of 1941, the first megawatt-class wind turbine was synchronized to a utility grid in Vermont. The Smith–Putnam wind turbine only ran for 1,100 hours before suffering a critical failure. The unit was not repaired, because of a shortage of materials during the war.
+
+The first utility grid-connected wind turbine to operate in the UK was built by John Brown & Company in 1951 in the Orkney Islands.[7][12]
+
+Despite these diverse developments, developments in fossil fuel systems almost entirely eliminated any wind turbine systems larger than supermicro size.  In the early 1970s, however, anti-nuclear protests in Denmark spurred artisan mechanics to develop microturbines of 22 kW.  Organizing owners into associations and co-operatives lead to the lobbying of the government and utilities and provided incentives for larger turbines throughout the 1980s and later.  Local activists in Germany, nascent turbine manufacturers in Spain, and large investors in the United States in the early 1990s then lobbied for policies that stimulated the industry in those countries.
+
+It has been argued that expanding use of wind power will lead to increasing geopolitical competition over critical materials for wind turbines such as rare earth elements neodymium, praseodymium, and dysprosium. But this perspective has been criticised for failing to recognise that most wind turbines do not use permanent magnets and for underestimating the power of economic incentives for expanded production of these minerals.[13]
+
+Wind Power Density (WPD) is a quantitative measure of  wind energy available at any location. It is the mean annual power available per square meter of swept area of a turbine, and is calculated for different heights above ground. Calculation of wind power density includes the effect of wind velocity and air density.[14]
+
+Wind turbines are classified by the wind speed they are designed for, from class I to class III, with A to C referring to the turbulence intensity of the wind.[15]
+
+Conservation of mass requires that the amount of air entering and exiting a turbine must be equal. Accordingly, Betz's law gives the maximal achievable extraction of wind power by a wind turbine as 16/27 (59.3%) of the rate at which the kinetic energy of the air arrives at the turbine.[16]
+
+The maximum theoretical power output of a wind machine is thus 16/27 times the rate at which kinetic energy of the air arrives at the effective disk area of the machine. If the effective area of the disk is A, and the wind velocity v, the maximum theoretical power output P is:
+
+where ρ is the air density.
+
+Wind-to-rotor efficiency (including rotor blade friction and drag) are among the factors affecting the final price of wind power.[17]
+Further inefficiencies, such as gearbox losses, generator and converter losses, reduce the power delivered by a wind turbine. To protect components from undue wear, extracted power is held constant above the rated operating speed as theoretical power increases at the cube of wind speed, further reducing theoretical efficiency. In 2001, commercial utility-connected turbines delivered 75% to 80% of the Betz limit of power extractable from the wind, at rated operating speed.[18][19][needs update]
+
+Efficiency can decrease slightly over time, one of the main reasons being dust and insect carcasses on the blades which alters the aerodynamic profile and essentially reduces the lift to drag ratio of the airfoil. Analysis of 3128 wind turbines older than 10 years in Denmark showed that half of the turbines had no decrease, while the other half saw a production decrease of 1.2% per year.[20]  Ice accretion on turbine blades has also been found to greatly reduce the efficiency of wind turbines, which is a common challenge in cold climates where in-cloud icing and freezing rain events occur.[21] Vertical turbine designs have much lower efficiency than standard horizontal designs.[22]
+
+In general, more stable and constant weather conditions (most notably wind speed) result in an average of 15% greater efficiency than that of a wind turbine in unstable weather conditions, thus allowing up to a 7% increase in wind speed under stable conditions. This is due to a faster recovery wake and greater flow entrainment that occur in conditions of higher atmospheric stability. However, wind turbine wakes have been found to recover faster under unstable atmospheric conditions as opposed to a stable environment.[23]
+
+Different materials have been found to have varying effects on the efficiency of wind turbines. In an Ege University experiment, three wind turbines (Each with three blades with diameters of one meter) were constructed with blades made of different materials: A glass and glass/carbon epoxy, glass/carbon, and glass/polyester. When tested, the results showed that the materials with higher overall masses had a greater friction moment and thus a lower power coefficient.[24]
+
+Wind turbines can rotate about either a horizontal or a vertical axis, the former being both older and more common.[25] They can also include blades, or be bladeless.[26] Vertical designs produce less power and are less common.[27]
+
+Large three-bladed horizontal-axis wind turbines (HAWT) with the blades upwind of the tower produce the overwhelming majority of wind power in the world today. These turbines have the main rotor shaft and electrical generator at the top of a tower, and must be pointed into the wind. Small turbines are pointed by a simple wind vane, while large turbines generally use a wind sensor coupled with a yaw system. Most have a gearbox, which turns the slow rotation of the blades into a quicker rotation that is more suitable to drive an electrical generator.[28] Some turbines use a different type of generator suited to slower rotational speed input. These don't need a gearbox and are called direct-drive, meaning they couple the rotor directly to the generator with no gearbox in between. While permanent magnet direct-drive generators can be more costly due to the rare earth materials required, these gearless turbines are sometimes preferred over gearbox generators because they "eliminate the gear-speed increaser, which is susceptible to significant accumulated fatigue torque loading, related reliability issues, and maintenance costs."[29] There is also the pseudo direct drive mechanism, which has some advantages over the permanent magnet direct drive mechanism.[30][31]
+
+Most horizontal axis turbines have their rotors upwind of the supporting tower. Downwind machines have been built, because they don't need an additional mechanism for keeping them in line with the wind. In high winds, the blades can also be allowed to bend, which reduces their swept area and thus their wind resistance. Despite these advantages, upwind designs are preferred, because the change in loading from the wind as each blade passes behind the supporting tower can cause damage to the turbine.
+
+Turbines used in wind farms for commercial production of electric power are usually three-bladed. These have low torque ripple, which contributes to good reliability. The blades are usually colored white for daytime visibility by aircraft and range in length from 20 to 80 meters (66 to 262 ft). The size and height of turbines increase year by year. Offshore wind turbines are built up to 8 MW today and have a blade length up to 80 meters (260 ft). Designs with 10 to 12 MW are in preparation.[32] Usual multi megawatt turbines have tubular steel towers with a height of 70 m to 120 m and in extremes up to 160 m.
+
+Vertical-axis wind turbines (or VAWTs) have the main rotor shaft arranged vertically. One advantage of this arrangement is that the turbine does not need to be pointed into the wind to be effective, which is an advantage on a site where the wind direction is highly variable. It is also an advantage when the turbine is integrated into a building because it is inherently less steerable.  Also, the generator and gearbox can be placed near the ground, using a direct drive from the rotor assembly to the ground-based gearbox, improving accessibility for maintenance. However, these designs produce much less energy averaged over time, which is a major drawback.[27][33]
+
+The key disadvantages include the relatively low rotational speed with the consequential higher torque and hence higher cost of the drive train, the inherently lower power coefficient, the 360-degree rotation of the aerofoil within the wind flow during each cycle and hence the highly dynamic loading on the blade, the pulsating torque generated by some rotor designs on the drive train, and the difficulty of modelling the wind flow accurately and hence the challenges of analysing and designing the rotor prior to fabricating a prototype.[34]
+
+When a turbine is mounted on a rooftop the building generally redirects wind over the roof and this can double the wind speed at the turbine. If the height of a rooftop mounted turbine tower is approximately 50% of the building height it is near the optimum for maximum wind energy and minimum wind turbulence. While wind speeds within the built environment are generally much lower than at exposed rural sites,[35][36] noise may be a concern and an existing structure may not adequately resist the additional stress.
+
+Subtypes of the vertical axis design include:
+
+"Eggbeater" turbines, or Darrieus turbines, were named after the French inventor, Georges Darrieus.[37] They have good efficiency, but produce large torque ripple and cyclical stress on the tower, which contributes to poor reliability. They also generally require some external power source, or an additional Savonius rotor to start turning, because the starting torque is very low. The torque ripple is reduced by using three or more blades, which results in greater solidity of the rotor. Solidity is measured by blade area divided by the rotor area. Newer Darrieus type turbines are not held up by guy-wires but have an external superstructure connected to the top bearing.[38]
+
+A subtype of Darrieus turbine with straight, as opposed to curved, blades. The cycloturbine variety has variable pitch to reduce the torque pulsation and is self-starting.[39] The advantages of variable pitch are: high starting torque; a wide, relatively flat torque curve; a higher coefficient of performance; more efficient operation in turbulent winds; and a lower blade speed ratio which lowers blade bending stresses. Straight, V, or curved blades may be used.[40]
+
+These are drag-type devices with two (or more) scoops that are used in anemometers, Flettner vents (commonly seen on bus and van roofs), and in some high-reliability low-efficiency power turbines. They are always self-starting if there are at least three scoops.
+
+Twisted Savonius is a modified savonius, with long helical scoops to provide smooth torque. This is often used as a rooftop wind turbine and has even been adapted for ships.[41]
+
+The parallel turbine is similar to the crossflow fan or centrifugal fan. It uses the ground effect. Vertical axis turbines of this type have been tried for many years: a unit producing 10 kW was built by Israeli wind pioneer Bruce Brill in the 1980s.[42][unreliable source?]
+
+Wind turbine design is a careful balance of cost, energy output, and fatigue life.
+
+Wind turbines convert wind energy to electrical energy for distribution. Conventional horizontal axis turbines can be divided into three components:
+
+A 1.5 (MW) wind turbine of a type frequently seen in the United States has a tower 80 meters (260 ft) high. The rotor assembly (blades and hub) weighs 22,000 kilograms (48,000 lb). The nacelle, which contains the generator, weighs 52,000 kilograms (115,000 lb). The concrete base for the tower is constructed using 26,000 kilograms (58,000 lb) reinforcing steel and contains 190 cubic meters (250 cu yd) of concrete. The base is 15 meters (50 ft) in diameter and 2.4 meters (8 ft) thick near the center.[48]
+
+Due to data transmission problems, structural health monitoring of wind turbines is usually performed using several accelerometers and strain gages attached to the nacelle to monitor the gearbox and equipment. Currently, digital image correlation and stereophotogrammetry are used to measure dynamics of wind turbine blades. These methods usually measure displacement and strain to identify location of defects. Dynamic characteristics of non-rotating wind turbines have been measured using digital image correlation and photogrammetry.[49] Three dimensional point tracking has also been used to measure rotating dynamics of wind turbines.[50]
+
+Wind turbine rotor blades are being made longer to increase efficiency. This requires them to be stiff, strong, light and resistant to fatigue.[51] Materials with these properties are composites such as polyester and epoxy, while glass fiber and carbon fiber have been used for the reinforcing.[52] Construction may use manual layup or injection molding.
+
+Companies seek ways to draw greater efficiency from their designs. A predominant way has been to increase blade length and thus rotor diameter. Retrofitting existing turbines with larger blades reduces the work and risks of redesigning the system. The current longest blade is 88.4 m (from LM Wind Power), but by 2021 offshore turbines are expected to be 10-MW with 100 m blades. Longer blades need to be stiffer to avoid deflection, which requires materials with higher stiffness-to-weight ratio. Because the blades need to function over a 100 million load cycles over a period of 20–25 years, the fatigue of the blade materials is also critical.
+
+Materials commonly used in wind turbine blades are described below.
+
+The stiffness of composites is determined by the stiffness of fibers and their volume content. Typically, E-glass fibers are used as main reinforcement in the composites. Typically, the glass/epoxy composites for wind turbine blades contain up to 75% glass by weight. This increases the stiffness, tensile and compression strength. A promising composite material is glass fiber with modified compositions like S-glass, R-glass etc. Other glass fibers developed by Owens Corning are ECRGLAS, Advantex and WindStrand.[53]
+
+Carbon fiber has more tensile strength, higher stiffness and lower density than glass fiber. An ideal candidate for these properties is the spar cap, a structural element of a blade which experiences high tensile loading.[52] A 100-m glass fiber blade could weigh up to 50 metric tons, while using carbon fiber in the spar saves 20% to 30% weight, about 15 metric tons.[54] However, because carbon fiber is ten times more expensive, glass fiber is still dominant.
+
+Instead of making wind turbine blade reinforcements from pure glass or pure carbon, hybrid designs trade weight for cost. For example, for an 8 m blade, a full replacement by carbon fiber would save 80% of weight but increase costs by 150%, while a 30% replacement would save 50% of weight and increase costs by 90%. Hybrid reinforcement materials include E-glass/carbon, E-glass/aramid. The current longest blade by LM Wind Power is made of carbon/glass hybrid composites. More research is needed about the optimal composition of materials [55]
+
+Additions of small amount (0.5 weight %) of nanoreinforcement (carbon nanotubes or nanoclay) in the polymer matrix of composites, fiber sizing or interlaminar layers can improve fatigue resistance, shear or compressive strength, and fracture toughness of the composites by 30% to 80%. Research has also shown that incorporating small amounts of carbon nanotubes (CNT) can increase the lifetime up to 1500%.
+
+As of 2019[update], a wind turbine may cost around $1 million per megawatt.[56]
+
+For the wind turbine blades, while the material cost is much higher for hybrid glass/carbon fiber blades than all-glass fiber blades, labor costs can be lower.  Using carbon fiber allows simpler designs that use less raw material.  The chief manufacturing process in blade fabrication is the layering of plies.  Thinner blades allow reducing the number of layers and so the labor, and in some cases, equate to the cost of labor for glass fiber blades.[57]
+
+Wind turbine parts other than the rotor blades (including the rotor hub, gearbox, frame, and tower) are largely made of steel. Smaller turbines (as well as megawatt-scale Enercon turbines) have begun using aluminum alloys for these components to make turbines lighter and more efficient. This trend may grow if fatigue and strength properties can be improved. 
+Pre-stressed concrete has been increasingly used for the material of the tower, but still requires much reinforcing steel to meet the strength requirement of the turbine. Additionally, step-up gearboxes are being increasingly replaced with variable speed generators, which requires magnetic materials.[51] In particular, this would require an greater supply of the rare earth metal neodymium.
+
+Modern turbines use a couple of tons of copper for generators, cables and such.[58] As of 2018[update], global production of wind turbines use 450,000 tonnes of copper per year.[59]
+
+A study of the material consumption trends and requirements for wind energy in Europe found that bigger turbines have a higher consumption of precious metals but lower material input per kW generated. The current material consumption and stock was compared to input materials for various onshore system sizes. In all EU countries the estimates for 2020 doubled the values consumed in 2009. These countries would need to expand their resources to meet the estimated demand for 2020. For example, currently the EU has 3% of world supply of fluorspar and it requires 14% by 2020. Globally, the main exporting countries are South Africa, Mexico and China.  This is similar with other critical and valuable materials required for energy systems such as magnesium, silver and indium. The levels of recycling of these materials are very low and focusing on that could alleviate supply. Because most of these valuable materials are also used  in other emerging technologies, like light emitting diodes (LEDs), photo voltaics (PVs) and liquid crystal displays (LCDs), their demand is expected to grow.[60]
+
+A study by the United States Geological Survey estimated resources required to fulfill the US commitment to supplying 20% of its electricity from wind power by 2030. It did not consider requirements for small turbines or offshore turbines because those were not common in 2008 when the study was done. Common materials such as cast iron, steel and concrete would increase by 2%–3% compared to 2008. Between 110,000 and 115,000 metric tons of fiber glass would be required per year, a 14% increase. Rare metal use would not increase much compared to available supply, however rare metals that are also used for other technologies such as batteries which are increasing its global demand need to be taken into account. Land required would be 50,000 square kilometers onshore and 11,000 offshore. This would not be a problem in the US due to its vast area and because the same land can be used for farming. A greater challenge would be the variability and transmission to areas of high demand.[61]
+
+Permanent magnets for wind turbine generators contain rare metals such as neodymium (Nd), praseodymium (Pr), Terbium (Tb) and dysprosium (Dy). Systems that use magnetic direct drive turbines require greater amounts of rare metals. Therefore, an increase in wind turbine manufacture would increase the demand for these resources. By 2035, the demand for Nd is estimated to increase by 4,000 to 18,000 tons and for Dy by 200 to 1200 tons. These values are a quarter to half of current production. However, these estimates are very uncertain because technologies are developing rapidly.[62]
+
+Reliance on rare earth minerals for components has risked expense and price volatility as China has been main producer of rare earth minerals (96% in 2009) and was reducing its export quotas.[63] However, in recent years other producers have increased production and China has increased export quotas, leading to a higher supply and lower cost, and a greater viability of large scale use of variable-speed generators.[64]
+
+Glass fiber is the most common material for reinforcement.  Its demand has grown due to growth in construction, transportation and wind turbines. Its global market might reach US$17.4 billion by 2024, compared to US$8.5 billion in 2014. In 2014, Asia Pacific produced more than 45% of the market; now China is the largest producer. The industry receives subsidies from the Chinese government allowing it to export cheaper to the US and Europe. However, price wars have led to anti-dumping measures such as tariffs on Chinese glass fiber.[65]
+
+Interest in recycling blades varies in different markets and depends on the waste legislation and local economics. A challenge in recycling blades is related to the composite material, which is made of a thermosetting matrix and glass fibers or a combination of glass and carbon fibers. Thermosetting matrix cannot be remolded to form new composites. So the options are either to send the blade to landfill, to reuse the blade and the composite material elements found in the blade, or to transform the composite material into a new source of material. In Germany, wind turbine blades are commercially recycled as part of an alternative fuel mix for a cement factory. In the USA the town of Casper, Wyoming has buried 1,000 non-recyclable blades in its landfill site, earning $675,000 for the town. It pointed out that wind farm waste is less toxic than other garbage. Wind turbine blades represent a “vanishingly small fraction” of overall waste in the US, according to the American Wind Energy Association.[66]
+
+A few localities have exploited the attention-getting nature of wind turbines by placing them on public display, either with visitor centers around their bases, or with viewing areas farther away.[67] The wind turbines are generally of conventional horizontal-axis, three-bladed design, and generate power to feed electrical grids, but they also serve the unconventional roles of technology demonstration, public relations, and education.
+
+Small wind turbines may be used for a variety of applications including on- or off-grid residences, telecom towers, offshore platforms, rural schools and clinics, remote monitoring and other purposes that require energy where there is no electric grid, or where the grid is unstable. Small wind turbines may be as small as a fifty-watt generator for boat or caravan use. Hybrid solar and wind powered units are increasingly being used for traffic signage, particularly in rural locations, as they avoid the need to lay long cables from the nearest mains connection point.[68] The U.S. Department of Energy's National Renewable Energy Laboratory (NREL) defines small wind turbines as those smaller than or equal to 100 kilowatts.[69] Small units often have direct drive generators, direct current output, aeroelastic blades, lifetime bearings and use a vane to point into the wind.
+
+Larger, more costly turbines generally have geared power trains, alternating current output, and flaps, and are actively pointed into the wind. Direct drive generators and aeroelastic blades for large wind turbines are being researched.
+
+On most horizontal wind turbine farms, a spacing of about 6–10 times the rotor diameter is often upheld. However, for large wind farms distances of about 15 rotor diameters should be more economical, taking into account typical wind turbine and land costs. This conclusion has been reached by research[70] conducted by Charles Meneveau of Johns Hopkins University[71] and Johan Meyers of Leuven University in Belgium, based on computer simulations[72] that take into account the detailed interactions among wind turbines (wakes) as well as with the entire turbulent atmospheric boundary layer.
+
+Recent research by John Dabiri of Caltech suggests that vertical wind turbines may be placed much more closely together so long as an alternating pattern of rotation is created allowing blades of neighbouring turbines to move in the same direction as they approach one another.[73]
+
+Wind turbines need regular maintenance to stay reliable and available. In the best case turbines are available to generate energy 98% of the time.[74][75]
+
+Modern turbines usually have a small onboard crane for hoisting maintenance tools and minor components. However, large, heavy components like generator, gearbox, blades, and so on are rarely replaced, and a heavy lift external crane is needed in those cases. If the turbine has a difficult access road, a containerized crane can be lifted up by the internal crane to provide heavier lifting.[76]
+
+Installation of new wind turbines can be controversial. An alternative is repowering, where existing wind turbines are replaced with bigger, more powerful ones, sometimes in smaller numbers while keeping or increasing capacity.
+
+Older turbines were in some early cases not required to be removed when reaching the end of their life. Some still stand, waiting to be recycled or repowered.[77][78]
+
+A demolition industry develops to recycle offshore turbines at a cost of DKK 2–4 million per (MW), to be guaranteed by the owner.[79]
+
+Wind turbines produce electricity at between two and six cents per kilowatt hour, which is one of the lowest-priced renewable energy sources.[80][81] As technology needed for wind turbines continued to improve, the prices decreased as well. In addition, there is currently no competitive market for wind energy, because wind is a freely available natural resource, most of which is untapped.[80] The main cost of small wind turbines is the purchase and installation process, which averages between $48,000 and $65,000 per installation. The energy harvested from the turbine will offset the installation cost, as well as provide virtually free energy for years.[82]
+
+Wind turbines provide a clean energy source, use little water,[1] emitting no greenhouse gases and no waste products. Over 1,500 tons of carbon dioxide per year can be eliminated by using a one-megawatt turbine instead of one megawatt of energy from a fossil fuel.[83]
+
+Wind turbines can be very large, reaching over 140 m (460 ft) tall and with blades 55 m (180 ft) long,[84] and people have often complained about their visual impact.
+
+Environmental impact of wind power includes effect on wildlife, but can be mitigated if proper monitoring and mitigation strategies are implemented.[85]  Thousands of birds, including rare species, have been killed by the blades of wind turbines,[86] though wind turbines contribute relatively insignificantly to anthropogenic avian mortality. Wind farms and nuclear power stations are responsible for between 0.3 and 0.4 bird deaths per gigawatt-hour (GWh) of electricity while fossil fueled power stations are responsible for about 5.2 fatalities per GWh. In 2009, for every bird killed by a wind turbine in the US, nearly 500,000 were killed by cats and another 500,000 by buildings.[87] In comparison, conventional coal fired generators contribute significantly more to bird mortality, by incineration when caught in updrafts of smoke stacks and by poisoning with emissions byproducts (including particulates and heavy metals downwind of flue gases). Further, marine life is affected by water intakes of steam turbine cooling towers (heat exchangers) for nuclear and fossil fuel generators, by coal dust deposits in marine ecosystems (e.g. damaging Australia's Great Barrier Reef) and by water acidification from combustion monoxides.
+
+Energy harnessed by wind turbines is intermittent, and is not a "dispatchable" source of power; its availability is based on whether the wind is blowing, not whether electricity is needed. Turbines can be placed on ridges or bluffs to maximize the access of wind they have, but this also limits the locations where they can be placed.[80] In this way, wind energy is not a particularly reliable source of energy. However, it can form part of the energy mix, which also includes power from other sources. Notably, the relative available output from wind and solar sources is often inversely proportional (balancing)[citation needed]. Technology is also being developed to store excess energy, which can then make up for any deficits in supplies.
+
+See also List of most powerful wind turbines

en/1782.html.txt

@@ -0,0 +1,145 @@
+
+
+A wind turbine, or alternatively referred to as a wind energy converter, is a device that converts the wind's kinetic energy into electrical energy.
+
+Wind turbines are manufactured  in a wide range of vertical and horizontal axis. The smallest turbines are used for applications such as battery charging for auxiliary power for boats or caravans or to power traffic warning signs. Larger turbines can be used for making contributions to a domestic power supply while selling unused power back to the utility supplier via the electrical grid. Arrays of large turbines, known as wind farms, are becoming an increasingly important source of intermittent  renewable energy and are used by many countries as part of a strategy to reduce their reliance on fossil fuels. One assessment claimed that, as of 2009[update], wind had the "lowest relative greenhouse gas emissions, the least water consumption demands and... the most favourable social impacts" compared to photovoltaic, hydro, geothermal, coal and gas.[1]
+
+The windwheel of Hero of Alexandria (10 AD – 70 AD) marks one of the first recorded instances of wind powering a machine in history.[2][3] However, the first known practical wind power plants were built in Sistan, an Eastern province of Persia (now Iran), from the 7th century. These "Panemone" were vertical axle windmills, which had long vertical drive shafts with rectangular blades.[4] Made of six to twelve sails covered in reed matting or cloth material, these windmills were used to grind grain or draw up water, and were used in the gristmilling and sugarcane industries.[5]
+
+Wind power first appeared in Europe during the Middle Ages. The first historical records of their use in England date to the 11th or 12th centuries, there are reports of German crusaders taking their windmill-making skills to Syria around 1190.[6] By the 14th century, Dutch windmills were in use to drain areas of the Rhine delta. Advanced wind turbines were described by Croatian inventor Fausto Veranzio. In his book Machinae Novae (1595) he described vertical axis wind turbines with curved or V-shaped blades.
+
+The first electricity-generating wind turbine was a battery charging machine installed in July 1887 by Scottish academic James Blyth to light his holiday home in Marykirk, Scotland.[7] Some months later American inventor Charles F. Brush was able to build the first automatically operated wind turbine after consulting local University professors and colleagues Jacob S. Gibbs and Brinsley Coleberd and successfully getting the blueprints peer-reviewed for electricity production in Cleveland, Ohio.[7] Although Blyth's turbine was considered uneconomical in the United Kingdom,[7] electricity generation by wind turbines was more cost effective in countries with widely scattered populations.[6]
+
+In Denmark by 1900, there were about 2500 windmills for mechanical loads such as pumps and mills, producing an estimated combined peak power of about 30 MW. The largest machines were on 24-meter (79 ft) towers with four-bladed 23-meter (75 ft) diameter rotors. By 1908, there were 72 wind-driven electric generators operating in the United States from 5 kW to 25 kW. Around the time of World War I, American windmill makers were producing 100,000 farm windmills each year, mostly for water-pumping.[9]
+
+By the 1930s, wind generators for electricity were common on farms, mostly in the United States where distribution systems had not yet been installed. In this period, high-tensile steel was cheap, and the generators were placed atop prefabricated open steel lattice towers.
+
+A forerunner of modern horizontal-axis wind generators was in service at Yalta, USSR in 1931. This was a 100 kW generator on a 30-meter (98 ft) tower, connected to the local 6.3 kV distribution system. It was reported to have an annual capacity factor of 32 percent, not much different from current wind machines.[10][11]
+
+In the autumn of 1941, the first megawatt-class wind turbine was synchronized to a utility grid in Vermont. The Smith–Putnam wind turbine only ran for 1,100 hours before suffering a critical failure. The unit was not repaired, because of a shortage of materials during the war.
+
+The first utility grid-connected wind turbine to operate in the UK was built by John Brown & Company in 1951 in the Orkney Islands.[7][12]
+
+Despite these diverse developments, developments in fossil fuel systems almost entirely eliminated any wind turbine systems larger than supermicro size.  In the early 1970s, however, anti-nuclear protests in Denmark spurred artisan mechanics to develop microturbines of 22 kW.  Organizing owners into associations and co-operatives lead to the lobbying of the government and utilities and provided incentives for larger turbines throughout the 1980s and later.  Local activists in Germany, nascent turbine manufacturers in Spain, and large investors in the United States in the early 1990s then lobbied for policies that stimulated the industry in those countries.
+
+It has been argued that expanding use of wind power will lead to increasing geopolitical competition over critical materials for wind turbines such as rare earth elements neodymium, praseodymium, and dysprosium. But this perspective has been criticised for failing to recognise that most wind turbines do not use permanent magnets and for underestimating the power of economic incentives for expanded production of these minerals.[13]
+
+Wind Power Density (WPD) is a quantitative measure of  wind energy available at any location. It is the mean annual power available per square meter of swept area of a turbine, and is calculated for different heights above ground. Calculation of wind power density includes the effect of wind velocity and air density.[14]
+
+Wind turbines are classified by the wind speed they are designed for, from class I to class III, with A to C referring to the turbulence intensity of the wind.[15]
+
+Conservation of mass requires that the amount of air entering and exiting a turbine must be equal. Accordingly, Betz's law gives the maximal achievable extraction of wind power by a wind turbine as 16/27 (59.3%) of the rate at which the kinetic energy of the air arrives at the turbine.[16]
+
+The maximum theoretical power output of a wind machine is thus 16/27 times the rate at which kinetic energy of the air arrives at the effective disk area of the machine. If the effective area of the disk is A, and the wind velocity v, the maximum theoretical power output P is:
+
+where ρ is the air density.
+
+Wind-to-rotor efficiency (including rotor blade friction and drag) are among the factors affecting the final price of wind power.[17]
+Further inefficiencies, such as gearbox losses, generator and converter losses, reduce the power delivered by a wind turbine. To protect components from undue wear, extracted power is held constant above the rated operating speed as theoretical power increases at the cube of wind speed, further reducing theoretical efficiency. In 2001, commercial utility-connected turbines delivered 75% to 80% of the Betz limit of power extractable from the wind, at rated operating speed.[18][19][needs update]
+
+Efficiency can decrease slightly over time, one of the main reasons being dust and insect carcasses on the blades which alters the aerodynamic profile and essentially reduces the lift to drag ratio of the airfoil. Analysis of 3128 wind turbines older than 10 years in Denmark showed that half of the turbines had no decrease, while the other half saw a production decrease of 1.2% per year.[20]  Ice accretion on turbine blades has also been found to greatly reduce the efficiency of wind turbines, which is a common challenge in cold climates where in-cloud icing and freezing rain events occur.[21] Vertical turbine designs have much lower efficiency than standard horizontal designs.[22]
+
+In general, more stable and constant weather conditions (most notably wind speed) result in an average of 15% greater efficiency than that of a wind turbine in unstable weather conditions, thus allowing up to a 7% increase in wind speed under stable conditions. This is due to a faster recovery wake and greater flow entrainment that occur in conditions of higher atmospheric stability. However, wind turbine wakes have been found to recover faster under unstable atmospheric conditions as opposed to a stable environment.[23]
+
+Different materials have been found to have varying effects on the efficiency of wind turbines. In an Ege University experiment, three wind turbines (Each with three blades with diameters of one meter) were constructed with blades made of different materials: A glass and glass/carbon epoxy, glass/carbon, and glass/polyester. When tested, the results showed that the materials with higher overall masses had a greater friction moment and thus a lower power coefficient.[24]
+
+Wind turbines can rotate about either a horizontal or a vertical axis, the former being both older and more common.[25] They can also include blades, or be bladeless.[26] Vertical designs produce less power and are less common.[27]
+
+Large three-bladed horizontal-axis wind turbines (HAWT) with the blades upwind of the tower produce the overwhelming majority of wind power in the world today. These turbines have the main rotor shaft and electrical generator at the top of a tower, and must be pointed into the wind. Small turbines are pointed by a simple wind vane, while large turbines generally use a wind sensor coupled with a yaw system. Most have a gearbox, which turns the slow rotation of the blades into a quicker rotation that is more suitable to drive an electrical generator.[28] Some turbines use a different type of generator suited to slower rotational speed input. These don't need a gearbox and are called direct-drive, meaning they couple the rotor directly to the generator with no gearbox in between. While permanent magnet direct-drive generators can be more costly due to the rare earth materials required, these gearless turbines are sometimes preferred over gearbox generators because they "eliminate the gear-speed increaser, which is susceptible to significant accumulated fatigue torque loading, related reliability issues, and maintenance costs."[29] There is also the pseudo direct drive mechanism, which has some advantages over the permanent magnet direct drive mechanism.[30][31]
+
+Most horizontal axis turbines have their rotors upwind of the supporting tower. Downwind machines have been built, because they don't need an additional mechanism for keeping them in line with the wind. In high winds, the blades can also be allowed to bend, which reduces their swept area and thus their wind resistance. Despite these advantages, upwind designs are preferred, because the change in loading from the wind as each blade passes behind the supporting tower can cause damage to the turbine.
+
+Turbines used in wind farms for commercial production of electric power are usually three-bladed. These have low torque ripple, which contributes to good reliability. The blades are usually colored white for daytime visibility by aircraft and range in length from 20 to 80 meters (66 to 262 ft). The size and height of turbines increase year by year. Offshore wind turbines are built up to 8 MW today and have a blade length up to 80 meters (260 ft). Designs with 10 to 12 MW are in preparation.[32] Usual multi megawatt turbines have tubular steel towers with a height of 70 m to 120 m and in extremes up to 160 m.
+
+Vertical-axis wind turbines (or VAWTs) have the main rotor shaft arranged vertically. One advantage of this arrangement is that the turbine does not need to be pointed into the wind to be effective, which is an advantage on a site where the wind direction is highly variable. It is also an advantage when the turbine is integrated into a building because it is inherently less steerable.  Also, the generator and gearbox can be placed near the ground, using a direct drive from the rotor assembly to the ground-based gearbox, improving accessibility for maintenance. However, these designs produce much less energy averaged over time, which is a major drawback.[27][33]
+
+The key disadvantages include the relatively low rotational speed with the consequential higher torque and hence higher cost of the drive train, the inherently lower power coefficient, the 360-degree rotation of the aerofoil within the wind flow during each cycle and hence the highly dynamic loading on the blade, the pulsating torque generated by some rotor designs on the drive train, and the difficulty of modelling the wind flow accurately and hence the challenges of analysing and designing the rotor prior to fabricating a prototype.[34]
+
+When a turbine is mounted on a rooftop the building generally redirects wind over the roof and this can double the wind speed at the turbine. If the height of a rooftop mounted turbine tower is approximately 50% of the building height it is near the optimum for maximum wind energy and minimum wind turbulence. While wind speeds within the built environment are generally much lower than at exposed rural sites,[35][36] noise may be a concern and an existing structure may not adequately resist the additional stress.
+
+Subtypes of the vertical axis design include:
+
+"Eggbeater" turbines, or Darrieus turbines, were named after the French inventor, Georges Darrieus.[37] They have good efficiency, but produce large torque ripple and cyclical stress on the tower, which contributes to poor reliability. They also generally require some external power source, or an additional Savonius rotor to start turning, because the starting torque is very low. The torque ripple is reduced by using three or more blades, which results in greater solidity of the rotor. Solidity is measured by blade area divided by the rotor area. Newer Darrieus type turbines are not held up by guy-wires but have an external superstructure connected to the top bearing.[38]
+
+A subtype of Darrieus turbine with straight, as opposed to curved, blades. The cycloturbine variety has variable pitch to reduce the torque pulsation and is self-starting.[39] The advantages of variable pitch are: high starting torque; a wide, relatively flat torque curve; a higher coefficient of performance; more efficient operation in turbulent winds; and a lower blade speed ratio which lowers blade bending stresses. Straight, V, or curved blades may be used.[40]
+
+These are drag-type devices with two (or more) scoops that are used in anemometers, Flettner vents (commonly seen on bus and van roofs), and in some high-reliability low-efficiency power turbines. They are always self-starting if there are at least three scoops.
+
+Twisted Savonius is a modified savonius, with long helical scoops to provide smooth torque. This is often used as a rooftop wind turbine and has even been adapted for ships.[41]
+
+The parallel turbine is similar to the crossflow fan or centrifugal fan. It uses the ground effect. Vertical axis turbines of this type have been tried for many years: a unit producing 10 kW was built by Israeli wind pioneer Bruce Brill in the 1980s.[42][unreliable source?]
+
+Wind turbine design is a careful balance of cost, energy output, and fatigue life.
+
+Wind turbines convert wind energy to electrical energy for distribution. Conventional horizontal axis turbines can be divided into three components:
+
+A 1.5 (MW) wind turbine of a type frequently seen in the United States has a tower 80 meters (260 ft) high. The rotor assembly (blades and hub) weighs 22,000 kilograms (48,000 lb). The nacelle, which contains the generator, weighs 52,000 kilograms (115,000 lb). The concrete base for the tower is constructed using 26,000 kilograms (58,000 lb) reinforcing steel and contains 190 cubic meters (250 cu yd) of concrete. The base is 15 meters (50 ft) in diameter and 2.4 meters (8 ft) thick near the center.[48]
+
+Due to data transmission problems, structural health monitoring of wind turbines is usually performed using several accelerometers and strain gages attached to the nacelle to monitor the gearbox and equipment. Currently, digital image correlation and stereophotogrammetry are used to measure dynamics of wind turbine blades. These methods usually measure displacement and strain to identify location of defects. Dynamic characteristics of non-rotating wind turbines have been measured using digital image correlation and photogrammetry.[49] Three dimensional point tracking has also been used to measure rotating dynamics of wind turbines.[50]
+
+Wind turbine rotor blades are being made longer to increase efficiency. This requires them to be stiff, strong, light and resistant to fatigue.[51] Materials with these properties are composites such as polyester and epoxy, while glass fiber and carbon fiber have been used for the reinforcing.[52] Construction may use manual layup or injection molding.
+
+Companies seek ways to draw greater efficiency from their designs. A predominant way has been to increase blade length and thus rotor diameter. Retrofitting existing turbines with larger blades reduces the work and risks of redesigning the system. The current longest blade is 88.4 m (from LM Wind Power), but by 2021 offshore turbines are expected to be 10-MW with 100 m blades. Longer blades need to be stiffer to avoid deflection, which requires materials with higher stiffness-to-weight ratio. Because the blades need to function over a 100 million load cycles over a period of 20–25 years, the fatigue of the blade materials is also critical.
+
+Materials commonly used in wind turbine blades are described below.
+
+The stiffness of composites is determined by the stiffness of fibers and their volume content. Typically, E-glass fibers are used as main reinforcement in the composites. Typically, the glass/epoxy composites for wind turbine blades contain up to 75% glass by weight. This increases the stiffness, tensile and compression strength. A promising composite material is glass fiber with modified compositions like S-glass, R-glass etc. Other glass fibers developed by Owens Corning are ECRGLAS, Advantex and WindStrand.[53]
+
+Carbon fiber has more tensile strength, higher stiffness and lower density than glass fiber. An ideal candidate for these properties is the spar cap, a structural element of a blade which experiences high tensile loading.[52] A 100-m glass fiber blade could weigh up to 50 metric tons, while using carbon fiber in the spar saves 20% to 30% weight, about 15 metric tons.[54] However, because carbon fiber is ten times more expensive, glass fiber is still dominant.
+
+Instead of making wind turbine blade reinforcements from pure glass or pure carbon, hybrid designs trade weight for cost. For example, for an 8 m blade, a full replacement by carbon fiber would save 80% of weight but increase costs by 150%, while a 30% replacement would save 50% of weight and increase costs by 90%. Hybrid reinforcement materials include E-glass/carbon, E-glass/aramid. The current longest blade by LM Wind Power is made of carbon/glass hybrid composites. More research is needed about the optimal composition of materials [55]
+
+Additions of small amount (0.5 weight %) of nanoreinforcement (carbon nanotubes or nanoclay) in the polymer matrix of composites, fiber sizing or interlaminar layers can improve fatigue resistance, shear or compressive strength, and fracture toughness of the composites by 30% to 80%. Research has also shown that incorporating small amounts of carbon nanotubes (CNT) can increase the lifetime up to 1500%.
+
+As of 2019[update], a wind turbine may cost around $1 million per megawatt.[56]
+
+For the wind turbine blades, while the material cost is much higher for hybrid glass/carbon fiber blades than all-glass fiber blades, labor costs can be lower.  Using carbon fiber allows simpler designs that use less raw material.  The chief manufacturing process in blade fabrication is the layering of plies.  Thinner blades allow reducing the number of layers and so the labor, and in some cases, equate to the cost of labor for glass fiber blades.[57]
+
+Wind turbine parts other than the rotor blades (including the rotor hub, gearbox, frame, and tower) are largely made of steel. Smaller turbines (as well as megawatt-scale Enercon turbines) have begun using aluminum alloys for these components to make turbines lighter and more efficient. This trend may grow if fatigue and strength properties can be improved. 
+Pre-stressed concrete has been increasingly used for the material of the tower, but still requires much reinforcing steel to meet the strength requirement of the turbine. Additionally, step-up gearboxes are being increasingly replaced with variable speed generators, which requires magnetic materials.[51] In particular, this would require an greater supply of the rare earth metal neodymium.
+
+Modern turbines use a couple of tons of copper for generators, cables and such.[58] As of 2018[update], global production of wind turbines use 450,000 tonnes of copper per year.[59]
+
+A study of the material consumption trends and requirements for wind energy in Europe found that bigger turbines have a higher consumption of precious metals but lower material input per kW generated. The current material consumption and stock was compared to input materials for various onshore system sizes. In all EU countries the estimates for 2020 doubled the values consumed in 2009. These countries would need to expand their resources to meet the estimated demand for 2020. For example, currently the EU has 3% of world supply of fluorspar and it requires 14% by 2020. Globally, the main exporting countries are South Africa, Mexico and China.  This is similar with other critical and valuable materials required for energy systems such as magnesium, silver and indium. The levels of recycling of these materials are very low and focusing on that could alleviate supply. Because most of these valuable materials are also used  in other emerging technologies, like light emitting diodes (LEDs), photo voltaics (PVs) and liquid crystal displays (LCDs), their demand is expected to grow.[60]
+
+A study by the United States Geological Survey estimated resources required to fulfill the US commitment to supplying 20% of its electricity from wind power by 2030. It did not consider requirements for small turbines or offshore turbines because those were not common in 2008 when the study was done. Common materials such as cast iron, steel and concrete would increase by 2%–3% compared to 2008. Between 110,000 and 115,000 metric tons of fiber glass would be required per year, a 14% increase. Rare metal use would not increase much compared to available supply, however rare metals that are also used for other technologies such as batteries which are increasing its global demand need to be taken into account. Land required would be 50,000 square kilometers onshore and 11,000 offshore. This would not be a problem in the US due to its vast area and because the same land can be used for farming. A greater challenge would be the variability and transmission to areas of high demand.[61]
+
+Permanent magnets for wind turbine generators contain rare metals such as neodymium (Nd), praseodymium (Pr), Terbium (Tb) and dysprosium (Dy). Systems that use magnetic direct drive turbines require greater amounts of rare metals. Therefore, an increase in wind turbine manufacture would increase the demand for these resources. By 2035, the demand for Nd is estimated to increase by 4,000 to 18,000 tons and for Dy by 200 to 1200 tons. These values are a quarter to half of current production. However, these estimates are very uncertain because technologies are developing rapidly.[62]
+
+Reliance on rare earth minerals for components has risked expense and price volatility as China has been main producer of rare earth minerals (96% in 2009) and was reducing its export quotas.[63] However, in recent years other producers have increased production and China has increased export quotas, leading to a higher supply and lower cost, and a greater viability of large scale use of variable-speed generators.[64]
+
+Glass fiber is the most common material for reinforcement.  Its demand has grown due to growth in construction, transportation and wind turbines. Its global market might reach US$17.4 billion by 2024, compared to US$8.5 billion in 2014. In 2014, Asia Pacific produced more than 45% of the market; now China is the largest producer. The industry receives subsidies from the Chinese government allowing it to export cheaper to the US and Europe. However, price wars have led to anti-dumping measures such as tariffs on Chinese glass fiber.[65]
+
+Interest in recycling blades varies in different markets and depends on the waste legislation and local economics. A challenge in recycling blades is related to the composite material, which is made of a thermosetting matrix and glass fibers or a combination of glass and carbon fibers. Thermosetting matrix cannot be remolded to form new composites. So the options are either to send the blade to landfill, to reuse the blade and the composite material elements found in the blade, or to transform the composite material into a new source of material. In Germany, wind turbine blades are commercially recycled as part of an alternative fuel mix for a cement factory. In the USA the town of Casper, Wyoming has buried 1,000 non-recyclable blades in its landfill site, earning $675,000 for the town. It pointed out that wind farm waste is less toxic than other garbage. Wind turbine blades represent a “vanishingly small fraction” of overall waste in the US, according to the American Wind Energy Association.[66]
+
+A few localities have exploited the attention-getting nature of wind turbines by placing them on public display, either with visitor centers around their bases, or with viewing areas farther away.[67] The wind turbines are generally of conventional horizontal-axis, three-bladed design, and generate power to feed electrical grids, but they also serve the unconventional roles of technology demonstration, public relations, and education.
+
+Small wind turbines may be used for a variety of applications including on- or off-grid residences, telecom towers, offshore platforms, rural schools and clinics, remote monitoring and other purposes that require energy where there is no electric grid, or where the grid is unstable. Small wind turbines may be as small as a fifty-watt generator for boat or caravan use. Hybrid solar and wind powered units are increasingly being used for traffic signage, particularly in rural locations, as they avoid the need to lay long cables from the nearest mains connection point.[68] The U.S. Department of Energy's National Renewable Energy Laboratory (NREL) defines small wind turbines as those smaller than or equal to 100 kilowatts.[69] Small units often have direct drive generators, direct current output, aeroelastic blades, lifetime bearings and use a vane to point into the wind.
+
+Larger, more costly turbines generally have geared power trains, alternating current output, and flaps, and are actively pointed into the wind. Direct drive generators and aeroelastic blades for large wind turbines are being researched.
+
+On most horizontal wind turbine farms, a spacing of about 6–10 times the rotor diameter is often upheld. However, for large wind farms distances of about 15 rotor diameters should be more economical, taking into account typical wind turbine and land costs. This conclusion has been reached by research[70] conducted by Charles Meneveau of Johns Hopkins University[71] and Johan Meyers of Leuven University in Belgium, based on computer simulations[72] that take into account the detailed interactions among wind turbines (wakes) as well as with the entire turbulent atmospheric boundary layer.
+
+Recent research by John Dabiri of Caltech suggests that vertical wind turbines may be placed much more closely together so long as an alternating pattern of rotation is created allowing blades of neighbouring turbines to move in the same direction as they approach one another.[73]
+
+Wind turbines need regular maintenance to stay reliable and available. In the best case turbines are available to generate energy 98% of the time.[74][75]
+
+Modern turbines usually have a small onboard crane for hoisting maintenance tools and minor components. However, large, heavy components like generator, gearbox, blades, and so on are rarely replaced, and a heavy lift external crane is needed in those cases. If the turbine has a difficult access road, a containerized crane can be lifted up by the internal crane to provide heavier lifting.[76]
+
+Installation of new wind turbines can be controversial. An alternative is repowering, where existing wind turbines are replaced with bigger, more powerful ones, sometimes in smaller numbers while keeping or increasing capacity.
+
+Older turbines were in some early cases not required to be removed when reaching the end of their life. Some still stand, waiting to be recycled or repowered.[77][78]
+
+A demolition industry develops to recycle offshore turbines at a cost of DKK 2–4 million per (MW), to be guaranteed by the owner.[79]
+
+Wind turbines produce electricity at between two and six cents per kilowatt hour, which is one of the lowest-priced renewable energy sources.[80][81] As technology needed for wind turbines continued to improve, the prices decreased as well. In addition, there is currently no competitive market for wind energy, because wind is a freely available natural resource, most of which is untapped.[80] The main cost of small wind turbines is the purchase and installation process, which averages between $48,000 and $65,000 per installation. The energy harvested from the turbine will offset the installation cost, as well as provide virtually free energy for years.[82]
+
+Wind turbines provide a clean energy source, use little water,[1] emitting no greenhouse gases and no waste products. Over 1,500 tons of carbon dioxide per year can be eliminated by using a one-megawatt turbine instead of one megawatt of energy from a fossil fuel.[83]
+
+Wind turbines can be very large, reaching over 140 m (460 ft) tall and with blades 55 m (180 ft) long,[84] and people have often complained about their visual impact.
+
+Environmental impact of wind power includes effect on wildlife, but can be mitigated if proper monitoring and mitigation strategies are implemented.[85]  Thousands of birds, including rare species, have been killed by the blades of wind turbines,[86] though wind turbines contribute relatively insignificantly to anthropogenic avian mortality. Wind farms and nuclear power stations are responsible for between 0.3 and 0.4 bird deaths per gigawatt-hour (GWh) of electricity while fossil fueled power stations are responsible for about 5.2 fatalities per GWh. In 2009, for every bird killed by a wind turbine in the US, nearly 500,000 were killed by cats and another 500,000 by buildings.[87] In comparison, conventional coal fired generators contribute significantly more to bird mortality, by incineration when caught in updrafts of smoke stacks and by poisoning with emissions byproducts (including particulates and heavy metals downwind of flue gases). Further, marine life is affected by water intakes of steam turbine cooling towers (heat exchangers) for nuclear and fossil fuel generators, by coal dust deposits in marine ecosystems (e.g. damaging Australia's Great Barrier Reef) and by water acidification from combustion monoxides.
+
+Energy harnessed by wind turbines is intermittent, and is not a "dispatchable" source of power; its availability is based on whether the wind is blowing, not whether electricity is needed. Turbines can be placed on ridges or bluffs to maximize the access of wind they have, but this also limits the locations where they can be placed.[80] In this way, wind energy is not a particularly reliable source of energy. However, it can form part of the energy mix, which also includes power from other sources. Notably, the relative available output from wind and solar sources is often inversely proportional (balancing)[citation needed]. Technology is also being developed to store excess energy, which can then make up for any deficits in supplies.
+
+See also List of most powerful wind turbines

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+The orca or killer whale (Orcinus orca) is a toothed whale belonging to the oceanic dolphin family, of which it is the largest member. Killer whales have a diverse diet, although individual populations often specialize in particular types of prey. Some feed exclusively on fish, while others hunt marine mammals such as seals and other species of dolphin. They have been known to attack baleen whale calves, and even adult whales. Killer whales are apex predators, as no animal preys on them.  A cosmopolitan species, they can be found in each of the world's oceans in a variety of marine environments, from Arctic and Antarctic regions to tropical seas, absent only from the Baltic and Black seas, and some areas of the Arctic Ocean.
+
+Killer whales are highly social; some populations are composed of matrilineal family groups (pods) which are the most stable of any animal species. Their sophisticated hunting techniques and vocal behaviours, which are often specific to a particular group and passed across generations, have been described as manifestations of animal culture.
+
+The International Union for Conservation of Nature assesses the orca's conservation status as data deficient because of the likelihood that two or more killer whale types are separate species. Some local populations are considered threatened or endangered due to prey depletion, habitat loss, pollution (by PCBs), capture for marine mammal parks, and conflicts with human fisheries. In late 2005, the southern resident killer whales, which swim in British Columbia and Washington state waters, were placed on the U.S. Endangered Species list.
+
+Wild killer whales are not considered a threat to humans and no fatal attack on humans has ever been documented, but there have been cases of captive orcas killing or injuring their handlers at marine theme parks. Killer whales feature strongly in the mythologies of indigenous cultures, and their reputation in different cultures ranges from being the souls of humans to merciless killers.
+
+Orcinus orca is the only recognized extant species in the genus Orcinus, and one of many animal species originally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae.[6] Konrad Gessner wrote the first scientific description of a killer whale in his Piscium & aquatilium animantium natura of 1558, part of the larger Historia animalium, based on examination of a dead stranded animal in the Bay of Greifswald that had attracted a great deal of local interest.[7]
+
+The killer whale is one of 35 species in the oceanic dolphin family, which first appeared about 11 million years ago. The killer whale lineage probably branched off shortly thereafter.[8] Although it has morphological similarities with the false killer whale, the pygmy killer whale and the pilot whales, a study of cytochrome b gene sequences by Richard LeDuc indicated that its closest extant relatives are the snubfin dolphins of the genus Orcaella.[9] However, a more recent (2018) study places the orca as a sister taxon to the Lissodelphininae, a clade that includes Lagenorhynchus and Cephalorhynchus.[10]
+
+Although the term "orca" is increasingly used, English-speaking scientists most often use the traditional name "killer whale"[citation needed]. The genus name Orcinus means "of the kingdom of the dead",[11] or "belonging to Orcus".[12]
+Ancient Romans originally used orca (pl. orcae) for these animals, possibly borrowing Ancient Greek ὄρυξ (óryx), which referred (among other things) to a whale species. Since the 1960s, "orca" has steadily grown in popularity. The term "orca" is preferred by some as it avoids the negative connotations of "killer",[13] and because, being part of the family Delphinidae, the species is more closely related to other oceanic dolphins than to other whales.[14]
+
+They are sometimes referred to as "blackfish", a name also used for other whale species. "Grampus" is a former name for the species, but is now seldom used. This meaning of "grampus" should not be confused with the genus Grampus, whose only member is Risso's dolphin.[15]
+
+The three to five types of killer whales may be distinct enough to be considered different races,[16] subspecies, or possibly even species[17] (see Species problem). The IUCN reported in 2008, "The taxonomy of this genus is clearly in need of review, and it is likely that O. orca will be split into a number of different species or at least subspecies over the next few years."[3] Although large variation in the ecological distinctiveness of different killer whale groups complicate simple differentiation into types,[18] research off the west coast of Canada and the United States in the 1970s and 1980s identified the following three types:
+
+Transients and residents live in the same areas, but avoid each other.[31][32][33]
+
+Other populations have not been as well studied, although specialized fish and mammal eating killer whales have been distinguished elsewhere.[34] In addition, separate populations of "generalist" (fish- and mammal-eating) and "specialist" (mammal-eating) killer whales have been identified off northwestern Europe.[35][36] As with residents and transients, the lifestyle of these whales appears to reflect their diet; fish-eating killer whales in Alaska[37] and Norway[38] have resident-like social structures, while mammal-eating killer whales in Argentina and the Crozet Islands behave more like transients.[39]
+
+Three types have been documented in the Antarctic. Two dwarf species, named Orcinus nanus and Orcinus glacialis, were described during the 1980s by Soviet researchers, but most cetacean researchers are sceptical about their status, and linking these directly to the types described below is difficult.[17]
+
+Types B and C live close to the ice pack, and diatoms in these waters may be responsible for the yellowish colouring of both types.[17][46] Mitochondrial DNA sequences support the theory that these are recently diverged separate species.[47] More recently, complete mitochondrial sequencing indicates the two Antarctic groups that eat seals and fish should be recognized as distinct species, as should the North Pacific transients, leaving the others as subspecies pending additional data.[48] Advanced methods that sequenced the entire mitochondrial genome revealed systematic differences in DNA between different populations.[49] A 2019 study of Type D orcas also found them to be distinct from other populations and possibly even a unique species.[44]
+
+Mammal-eating killer whales in different regions were long thought likely to be closely related, but genetic testing has refuted this hypothesis.[50]
+
+There are seven identified ecotypes inhabiting isolated ecological niches. Of three orca ecotypes in the Antarctic, one preys on minke whales, the second on seals and penguins, and the third on fish. Another ecotype lives in the eastern North Atlantic, while the three Northeast Pacific ecotypes are labelled the transient, resident and offshore populations described above. Research has supported a proposal to reclassify the Antarctic seal- and fish-eating populations and the North Pacific transients as a distinct species, leaving the remaining ecotypes as subspecies. The first split in the orca population, between the North Pacific transients and the rest, occurred an estimated 700,000 years ago. Such a designation would mean that each new species becomes subject to separate conservation assessments.[49]
+
+A typical killer whale distinctively bears a black back, white chest and sides, and a white patch above and behind the eye. Calves are born with a yellowish or orange tint, which fades to white. It has a heavy and robust body[51] with a large dorsal fin up to 1.8 m (5 ft 11 in) tall.[52] Behind the fin, it has a dark grey "saddle patch" across the back. Antarctic killer whales may have pale grey to nearly white backs. Adult killer whales are very distinctive, seldom confused with any other sea creature.[53] When seen from a distance, juveniles can be confused with other cetacean species, such as the false killer whale or Risso's dolphin.[54]
+
+The killer whale's teeth are very strong, and its jaws exert a powerful grip; the upper teeth fall into the gaps between the lower teeth when the mouth is closed. The firm middle and back teeth hold prey in place, while the front teeth are inclined slightly forward and outward to protect them from powerful jerking movements.[55]
+
+Killer whales are the largest extant members of the dolphin family. Males typically range from 6 to 8 metres (20 to 26 ft) long and weigh in excess of 6 tonnes (5.9 long tons; 6.6 short tons). Females are smaller, generally ranging from 5 to 7 m (16 to 23 ft) and weighing about 3 to 4 tonnes (3.0 to 3.9 long tons; 3.3 to 4.4 short tons).[56]  Calves at birth weigh about 180 kg (400 lb) and are about 2.4 m (7.9 ft) long.[57][58] The skeleton of the killer whale is of the typical delphinid structure, but more robust.[59] Its integument, unlike that of most other dolphin species, is characterized by a well-developed dermal layer with a dense network of fascicles of collagen fibres.[60]
+
+Killer whale pectoral fins, analogous to forelimbs, are large and rounded, resembling paddles, with those of males significantly larger than those of females. Dorsal fins also exhibit sexual dimorphism, with those of males about 1.8 m (5.9 ft) high, more than twice the size of the female's, with the male's fin more like a tall, elongated isosceles triangle, whereas the female's is shorter and more curved.[61] Males and females also have different patterns of black and white skin in their genital areas.[62] In the skull, adult males have longer lower jaws than females, as well as larger occipital crests.[60]
+
+An individual killer whale can often be identified from its dorsal fin and saddle patch. Variations such as nicks, scratches, and tears on the dorsal fin and the pattern of white or grey in the saddle patch are unique. Published directories contain identifying photographs and names for hundreds of North Pacific animals. Photographic identification has enabled the local population of killer whales to be counted each year rather than estimated, and has enabled great insight into life cycles and social structures.[63]
+
+Occasionally a killer whale is white; they have been spotted in the northern Bering Sea and around St. Lawrence Island, and near the Russian coast.[64][65] In February 2008, a white killer whale was photographed 3.2 km (2.0 mi) off Kanaga Volcano in the Aleutian Islands.[64][65]  In 2010, the Far East Russia Orca Project (FEROP), co-founded and co-directed by Alexander M. Burdin and Erich Hoyt, filmed an adult male nicknamed Iceberg.[66][67]
+
+Killer whales have good eyesight above and below the water, excellent hearing, and a good sense of touch. They have exceptionally sophisticated echolocation abilities, detecting the location and characteristics of prey and other objects in the water by emitting clicks and listening for echoes,[68] as do other members of the dolphin family. The mean body temperature of the orca is 36 to 38 °C (97 to 100 °F).[69][70] Like most marine mammals, orcas have a layer of insulating blubber ranging from 7.6 to 10 cm (3.0 to 3.9 in) thick[69] beneath the skin. The pulse is about 60 heartbeats per minute when the orca is at the surface, dropping to 30 beats/min when submerged.[71]
+
+Killer whales are found in all oceans and most seas. Due to their enormous range, numbers, and density, relative distribution is difficult to estimate,[72] but they clearly prefer higher latitudes and coastal areas over pelagic environments.[73] Areas which serve as major study sites for the species include the coasts of Iceland, Norway, the Valdes Peninsula of Argentina, the Crozet Islands, New Zealand and parts of the west coast of North America, from California to Alaska.[74]
+
+Systematic surveys indicate the highest densities of killer whales (>0.40 individuals per 100 km2) in the northeast Atlantic around the Norwegian coast, in the north Pacific along the Aleutian Islands, the Gulf of Alaska and in the Southern Ocean off much of the coast of Antarctica.[72] They are considered "common" (0.20–0.40 individuals per 100 km2) in the eastern Pacific along the coasts of British Columbia, Washington and Oregon, in the North Atlantic Ocean around Iceland and the Faroe Islands. High densities have also been reported but not quantified in the western North Pacific around the Sea of Japan, Sea of Okhotsk, Kuril Islands, Kamchatka and the Commander Islands and in the Southern Hemisphere off southern Brazil and the tip of southern Africa. They are reported as seasonally common in the Canadian Arctic, including Baffin Bay between Greenland and Nunavut, as well as Tasmania and Macquarie Island.[72] Regularly occurring or distinct populations exist off Northwest Europe, California, Patagonia, the Crozet Islands, Marion Island, southern Australia and New Zealand.[36][72][75] The northwest Atlantic population of at least 67 individuals ranges from Labrador and Newfoundland to New England with sightings to Cape Cod and Long Island.[76]
+
+Information for offshore regions and warmer waters is more scarce, but widespread sightings indicate that the killer whale can survive in most water temperatures. They have been sighted, though more infrequently, in the Mediterranean, the Arabian Sea, the Gulf of Mexico, Banderas Bay on Mexico's west coast and the Caribbean.[72] Over 50 individual whales have been documented in the northern Indian Ocean, including two individuals that were sighted in the Persian Gulf in 2008 and off Sri Lanka in 2015.[77] Those orcas may occasionally enter the Red Sea through the Gulf of Aden.[78] The modern status of the species along coastal mainland China and its vicinity is unknown. Recorded sightings have been made from almost the entire shoreline.[79] A wide-ranging population is likely to exist in the central Pacific, with some sightings off Hawaii.[80][81] Distinct populations may also exist off the west coast of tropical Africa,[82] and Papua New Guinea.[83] In the Mediterranean, killer whales are considered "visitors", likely from the North Atlantic, and sightings become less frequent further east. However, a small year-round population is known to exist in the Strait of Gibraltar, mostly on the Atlantic side.[84][85] Killer whales also appear to regularly occur off the Galápagos Islands.[86]
+
+In the Antarctic, killer whales range up to the edge of the pack ice and are believed to venture into the denser pack ice, finding open leads much like beluga whales in the Arctic. However, killer whales are merely seasonal visitors to Arctic waters, and do not approach the pack ice in the summer. With the rapid Arctic sea ice decline in the Hudson Strait, their range now extends deep into the northwest Atlantic.[87] Occasionally, killer whales swim into freshwater rivers. They have been documented 100 mi (160 km) up the Columbia River in the United States.[88][89] They have also been found in the Fraser River in Canada and the Horikawa River in Japan.[88]
+
+Migration patterns are poorly understood. Each summer, the same individuals appear off the coasts of British Columbia and Washington. Despite decades of research, where these animals go for the rest of the year remains unknown. Transient pods have been sighted from southern Alaska to central California.[90]
+
+Worldwide population estimates are uncertain, but recent consensus suggests a minimum of 50,000 (2006).[91][3][30] Local estimates include roughly 25,000 in the Antarctic, 8,500 in the tropical Pacific, 2,250–2,700 off the cooler northeast Pacific and 500–1,500 off Norway.[92] Japan's Fisheries Agency estimated in the 2000s that 2,321 killer whales were in the seas around Japan.[93][94]
+
+Killer whales are apex predators, meaning that they themselves have no natural predators. They are sometimes called the wolves of the sea, because they hunt in groups like wolf packs.[95] Killer whales hunt varied prey including fish, cephalopods, mammals, seabirds, and sea turtles.[96] Different populations or ecotypes may specialize, and some can have a dramatic impact on prey species.[97] However, whales in tropical areas appear to have more generalized diets due to lower food productivity.[81][82] Killer whales spend most of their time at shallow depths,[98] but occasionally dive several hundred meters depending on their prey.[99][100]
+
+Fish-eating killer whales prey on around 30 species of fish. Some populations in the Norwegian and Greenland sea specialize in herring and follow that fish's autumnal migration to the Norwegian coast. Salmon account for 96% of northeast Pacific residents' diet, including 65% of large, fatty Chinook.[101] Chum salmon are also eaten, but smaller sockeye and pink salmon are not a significant food item.[102] Depletion of specific prey species in an area is, therefore, cause for concern for local populations, despite the high diversity of prey. On average, a killer whale eats 227 kilograms (500 lb) each day.[103] While salmon are usually hunted by an individual whale or a small group, herring are often caught using carousel feeding: the killer whales force the herring into a tight ball by releasing bursts of bubbles or flashing their white undersides. They then slap the ball with their tail flukes, stunning or killing up to 15 fish at a time, then eating them one by one. Carousel feeding has only been documented in the Norwegian killer whale population, as well as some oceanic dolphin species.[104]
+
+In New Zealand, sharks and rays appear to be important prey, including eagle rays, long-tail and short-tail stingrays, common threshers, smooth hammerheads, blue sharks, basking sharks, and shortfin makos.[105][106] With sharks, orcas may herd them to the surface and strike them with their tail flukes,[105] while bottom-dwelling rays are cornered, pinned to the ground and taken to the surface.[107] In other parts of the world, killer whales have preyed on broadnose sevengill sharks,[108] tiger sharks[109] and even small whale sharks.[110] Killer whales have also been recorded attacking and feeding on great white sharks,[28][111][112][113] and appear to target the liver.[28][112] Competition between killer whales and white sharks is probable in regions where their diets overlap.[114] The arrival of orcas in an area can cause white sharks to flee and forage elsewhere.[115]
+
+Killer whales are very sophisticated and effective predators of marine mammals. Thirty-two cetacean species have been recorded as prey, from observing orcas' feeding activity, examining the stomach contents of dead orcas, and seeing scars on the bodies of surviving prey animals. Groups even attack larger cetaceans such as minke whales, grey whales,[116] and, rarely, sperm whales or blue whales.[34][117][118][119] It has been hypothesized that predation by orcas on whale calves in high-productivity, high-latitude areas is the reason for great whale migrations during breeding season to low-productivity tropical waters where orcas are scarcer.[120]
+
+Hunting a large whale usually takes several hours. Killer whales generally attack young or weak animals; however, a group of five or more may attack a healthy adult. When hunting a young whale, a group chases it and its mother to exhaustion. Eventually, they separate the pair and surround the calf, drowning it by keeping it from surfacing. Pods of female sperm whales sometimes protect themselves by forming a protective circle around their calves with their flukes facing outwards, using them to repel the attackers.[121] Rarely, large killer whale pods can overwhelm even adult female sperm whales. Adult bull sperm whales, which are large, powerful and aggressive when threatened, and fully grown adult blue whales, which are possibly too large to overwhelm, are not believed to be prey for killer whales.[122]
+
+Prior to the advent of industrial whaling, great whales may have been the major food source for killer whales. The introduction of modern whaling techniques may have aided killer whales by the sound of exploding harpoons indicating availability of prey to scavenge, and compressed air inflation of whale carcasses causing them to float, thus exposing them to scavenging. However, the devastation of great whale populations by unfettered whaling has possibly reduced their availability for killer whales, and caused them to expand their consumption of smaller marine mammals, thus contributing to the decline of these as well.[120]
+
+Other marine mammal prey species include nearly 20 species of seal, sea lion and fur seal. Walruses and sea otters are less frequently taken. Often, to avoid injury, killer whales disable their prey before killing and eating it. This may involve throwing it in the air, slapping it with their tails, ramming it, or breaching and landing on it.[123] In the Aleutian Islands, a decline in sea otter populations in the 1990s was controversially attributed by some scientists to killer whale predation, although with no direct evidence.[124] The decline of sea otters followed a decline in harbour seal and Steller sea lion populations, the killer whale's preferred prey,[a][126] which in turn may be substitutes for their original prey, now decimated by industrial whaling.[127][128][129]
+
+In steeply banked beaches off Península Valdés, Argentina, and the Crozet Islands, killer whales feed on South American sea lions and southern elephant seals in shallow water, even beaching temporarily to grab prey before wriggling back to the sea. Beaching, usually fatal to cetaceans, is not an instinctive behaviour, and can require years of practice for the young.[130] Killer whales can then release the animal near juvenile whales, allowing the younger whales to practice the difficult capture technique on the now-weakened prey.[123][131] "Wave-hunting" killer whales spy-hop to locate Weddell seals, crabeater seals, leopard seals, and penguins resting on ice floes, and then swim in groups to create waves that wash over the floe. This washes the prey into the water, where other killer whales lie in wait.[49][132][133]
+
+Killer whales have also been observed preying on terrestrial mammals, such as deer swimming between islands off the northwest coast of North America.[125] Killer whale cannibalism has also been reported based on analysis of stomach contents, but this is likely to be the result of scavenging remains dumped by whalers.[134] One killer whale was also attacked by its companions after being shot.[34] Although resident killer whales have never been observed to eat other marine mammals, they occasionally harass and kill porpoises and seals for no apparent reason.[135]
+
+Killer whales in many areas may prey on cormorants and gulls.[136] A captive killer whale at Marineland of Canada discovered it could regurgitate fish onto the surface, attracting sea gulls, and then eat the birds. Four others then learned to copy the behaviour.[137]
+
+Day-to-day killer whale behaviour generally consists of foraging, travelling, resting and socializing. Killer whales frequently engage in surface behaviour such as breaching (jumping completely out of the water) and tail-slapping. These activities may have a variety of purposes, such as courtship, communication, dislodging parasites, or play. Spyhopping is a behaviour in which a whale holds its head above water to view its surroundings.[138] Resident killer whales swim alongside porpoises and other dolphins.[139]
+
+Killer whales are notable for their complex societies. Only elephants and higher primates live in comparably complex social structures.[140] Due to orcas' complex social bonds, many marine experts have concerns about how humane it is to keep them in captivity.[141]
+
+Resident killer whales in the eastern North Pacific live in particularly complex and stable social groups. Unlike any other known mammal social structure, resident whales live with their mothers for their entire lives. These family groups are based on matrilines consisting of the eldest female (matriarch) and her sons and daughters, and the descendants of her daughters, etc. The average size of a matriline is 5.5 animals.[142] Because females can reach age 90, as many as four generations travel together. These matrilineal groups are highly stable. Individuals separate for only a few hours at a time, to mate or forage. With one exception, a killer whale named Luna, no permanent separation of an individual from a resident matriline has been recorded.[142]
+
+Closely related matrilines form loose aggregations called pods, usually consisting of one to four matrilines. Unlike matrilines, pods may separate for weeks or months at a time.[142] DNA testing indicates resident males nearly always mate with females from other pods.[143] Clans, the next level of resident social structure, are composed of pods with similar dialects, and common but older maternal heritage. Clan ranges overlap, mingling pods from different clans.[142] The final association layer, perhaps more arbitrarily defined than the familial groupings, is called the community, and is defined as a set of clans that regularly commingle. Clans within a community do not share vocal patterns.[b]
+
+Transient pods are smaller than resident pods, typically consisting of an adult female and one or two of her offspring. Males typically maintain stronger relationships with their mothers than other females. These bonds can extend well into adulthood. Unlike residents, extended or permanent separation of transient offspring from natal matrilines is common, with juveniles and adults of both sexes participating. Some males become "rovers" and do not form long-term associations, occasionally joining groups that contain reproductive females.[144] As in resident clans, transient community members share an acoustic repertoire, although regional differences in vocalizations have been noted.[145]
+
+Like all cetaceans, killer whales depend heavily on underwater sound for orientation, feeding, and communication. They produce three categories of sounds: clicks, whistles, and pulsed calls. Clicks are believed to be used primarily for navigation and discriminating prey and other objects in the surrounding environment, but are also commonly heard during social interactions.[30]
+
+Northeast Pacific resident groups tend to be much more vocal than transient groups in the same waters.[146] Residents feed primarily on Chinook and chum, which are insensitive to killer whale calls (inferred from the audiogram of Atlantic salmon). In contrast, the marine mammal prey of transients hear whale calls well. Transients are typically silent.[146] They sometimes use a single click (called a cryptic click) rather than the long train of clicks observed in other populations. Residents are silent only when resting.
+
+All members of a resident pod use similar calls, known collectively as a dialect. Dialects are composed of specific numbers and types of discrete, repetitive calls. They are complex and stable over time.[147] Call patterns and structure are distinctive within matrilines.[148] Newborns produce calls similar to their mothers, but have a more limited repertoire.[145] Individuals likely learn their dialect through contact with pod members.[149] Family-specific calls have been observed more frequently in the days following a calf's birth, which may help the calf learn them.[150] Dialects are probably an important means of maintaining group identity and cohesiveness. Similarity in dialects likely reflects the degree of relatedness between pods, with variation growing over time.[151] When pods meet, dominant call types decrease and subset call types increase. The use of both call types is called biphonation. The increased subset call types may be the distinguishing factor between pods and inter-pod relations.[148]
+
+Dialects also distinguish types. Resident dialects contain seven to 17 (mean = 11) distinctive call types. All members of the North American west coast transient community express the same basic dialect, although minor regional variation in call types is evident. Preliminary research indicates offshore killer whales have group-specific dialects unlike those of residents and transients.[151]
+
+Norwegian and Icelandic herring-eating orcas appear to have different vocalizations for activities like hunting.[152] A population that live in McMurdo Sound, Antarctica have 28 complex burst-pulse and whistle calls.[153]
+
+Killer whales have the second-heaviest brains among marine mammals[154] (after sperm whales, which have the largest brain of any animal). They can be trained in captivity and are often described as intelligent,[155][156] although defining and measuring "intelligence" is difficult in a species whose environment and behavioural strategies are very different from those of humans.[156]
+
+Killer whales imitate others, and seem to deliberately teach skills to their kin. Off the Crozet Islands, mothers push their calves onto the beach, waiting to pull the youngster back if needed.[123][131]
+
+People who have interacted closely with killer whales offer numerous anecdotes demonstrating the whales' curiosity, playfulness, and ability to solve problems. Alaskan killer whales have not only learned how to steal fish from longlines, but have also overcome a variety of techniques designed to stop them, such as the use of unbaited lines as decoys.[157] Once, fishermen placed their boats several miles apart, taking turns retrieving small amounts of their catch, in the hope that the whales would not have enough time to move between boats to steal the catch as it was being retrieved. A researcher described what happened next:
+
+It worked really well for a while. Then the whales split into two groups. It didn't even take them an hour to figure it out. They were so thrilled when they figured out what was going on, that we were playing games. They were breaching by the boats.
+
+In other anecdotes, researchers describe incidents in which wild killer whales playfully tease humans by repeatedly moving objects the humans are trying to reach,[158] or suddenly start to toss around a chunk of ice after a human throws a snowball.[159]
+
+The killer whale's use of dialects and the passing of other learned behaviours from generation to generation have been described as a form of animal culture.[160]
+
+The complex and stable vocal and behavioural cultures of sympatric groups of killer whales (Orcinus orca) appear to have no parallel outside humans and represent an independent evolution of cultural faculties.[161]
+
+(Two species or populations are considered sympatric when they live in the same geographic area and thus regularly encounter one another.)
+
+Female killer whales begin to mature at around the age of 10 and reach peak fertility around 20,[162] experiencing periods of polyestrous cycling separated by non-cycling periods of three to 16 months. Females can often breed until age 40, followed by a rapid decrease in fertility.[162] As such, orcas are among the few animals that undergo menopause and live for decades after they have finished breeding.[163][164] The lifespans of wild females average 50 to 80 years.[165] Some are claimed to have lived substantially longer: Granny (J2) was estimated by some researchers to have been as old as 105 years at the time of her death, though a biopsy sample indicated her age as 65 to 80 years.[166][167][168] Orcas held in captivity tend to live less than those in the wild, although this is subject to scientific debate.[165][169][170]
+
+To avoid inbreeding, males mate with females from other pods. Gestation varies from 15 to 18 months. [171] Mothers usually calve a single offspring about once every five years. In resident pods, births occur at any time of year, although winter is the most common. Mortality is extremely high during the first seven months of life, when 37–50% of all calves die.[172] Weaning begins at about 12 months of age, and is complete by two years. According to observations in several regions, all male and female pod members participate in the care of the young.[140]
+
+Males sexually mature at the age of 15, but do not typically reproduce until age 21. Wild males live around 29 years on average, with a maximum of about 60 years.[166] One male, known as Old Tom, was reportedly spotted every winter between the 1840s and 1930 off New South Wales, Australia. This would have made him up to 90 years old. Examination of his teeth indicated he died around age 35,[173] but this method of age determination is now believed to be inaccurate for older animals.[174] One male known to researchers in the Pacific Northwest (identified as J1) was estimated to have been 59 years old when he died in 2010.[175] Killer whales are unique among cetaceans, as their caudal sections elongate with age, making their heads relatively shorter.[60]
+
+Infanticide, once thought to occur only in captive killer whales, was observed in wild populations by researchers off British Columbia on December 2, 2016. In this incident, an adult male killed the calf of a female within the same pod, with his mother also joining in the assault. It is theorized that the male killed the young calf in order to mate with its mother (something that occurs in other carnivore species), while the male's mother supported the breeding opportunity for her son.  The attack ended when the calf's mother struck and injured the attacking male. Such behaviour matches that of many smaller dolphin species, such as the bottlenose dolphin.[176]
+
+In 2008, the IUCN (International Union for Conservation of Nature) changed its assessment of the killer whale's conservation status from conservation dependent to data deficient, recognizing that one or more killer whale types may actually be separate, endangered species.[3] Depletion of prey species, pollution, large-scale oil spills, and habitat disturbance caused by noise and conflicts with boats are the most significant worldwide threats.[3] In January 2020, the first killer whale in England and Wales since 2001 was found dead with a large fragment of plastic in its stomach.[177]
+
+Like other animals at the highest trophic levels, the killer whale is particularly at risk of poisoning from bioaccumulation of toxins, including Polychlorinated biphenyls (PCBs).[178] European harbour seals have problems in reproductive and immune functions associated with high levels of PCBs and related contaminants, and a survey off the Washington coast found PCB levels in killer whales were higher than levels that had caused health problems in harbour seals.[178] Blubber samples in the Norwegian Arctic show higher levels of PCBs, pesticides and brominated flame-retardants than in polar bears. When food is scarce, killer whales metabolize blubber for energy, which increases pollutant concentrations in their blood.
+
+In the Pacific Northwest, wild salmon stocks, a main resident food source, have declined dramatically in recent years.[3]  In the Puget Sound region only 75 whales remain with few births over the last few years.[179] On the west coast of Alaska and the Aleutian Islands, seal and sea lion populations have also substantially declined.[180]
+
+In 2005, the United States government listed the southern resident community as an endangered population under the Endangered Species Act.[30] This community comprises three pods which live mostly in the Georgia and Haro Straits and Puget Sound in British Columbia and Washington. They do not breed outside of their community, which was once estimated at around 200 animals and later shrank to around 90.[181] In October 2008, the annual survey revealed seven were missing and presumed dead, reducing the count to 83.[182] This is potentially the largest decline in the population in the past 10 years. These deaths can be attributed to declines in Chinook salmon.[182]
+
+Scientist Ken Balcomb has extensively studied killer whales since 1976; he is the research biologist responsible for discovering U.S. Navy sonar may harm killer whales. He studied killer whales from the Center for Whale Research, located in Friday Harbor, Washington.[183] He was also able to study killer whales from "his home porch perched above Puget Sound, where the animals hunt and play in summer months".[183] In May 2003, Balcomb (along with other whale watchers near the Puget Sound coastline) noticed uncharacteristic behaviour displayed by the killer whales. The whales seemed "agitated and were moving haphazardly, attempting to lift their heads free of the water" to escape the sound of the sonars.[183] "Balcomb confirmed at the time that strange underwater pinging noises detected with underwater microphones were sonar. The sound originated from a U.S. Navy frigate 12 miles (19 kilometres) distant, Balcomb said."[183] The impact of sonar waves on killer whales is potentially life-threatening. Three years prior to Balcomb's discovery, research in the Bahamas showed 14 beaked whales washed up on the shore. These whales were beached on the day U.S. Navy destroyers were activated into sonar exercise.[183] Of the 14 whales beached, six of them died. These six dead whales were studied, and CAT scans of two of the whale heads showed hemorrhaging around the brain and the ears, which is consistent with decompression sickness.[183]
+
+Another conservation concern was made public in September 2008 when the Canadian government decided it was not necessary to enforce further protections (including the Species at Risk Act in place to protect endangered animals along their habitats) for killer whales aside from the laws already in place. In response to this decision, six environmental groups sued the federal government, claiming killer whales were facing many threats on the British Columbia Coast and the federal government did nothing to protect them from these threats.[184] A legal and scientific nonprofit organization, Ecojustice, led the lawsuit and represented the David Suzuki Foundation, Environmental Defence, Greenpeace Canada, International Fund for Animal Welfare, the Raincoast Conservation Foundation, and the Wilderness Committee.[184] Many scientists involved in this lawsuit, including Bill Wareham, a marine scientist with the David Suzuki Foundation, noted increased boat traffic, water toxic wastes, and low salmon population as major threats, putting approximately 87 killer whales[184] on the British Columbia Coast in danger.
+
+Underwater noise from shipping, drilling, and other human activities is a significant concern in some key killer whale habitats, including Johnstone Strait and Haro Strait.[185] In the mid-1990s, loud underwater noises from salmon farms were used to deter seals. Killer whales also avoided the surrounding waters.[186] High-intensity sonar used by the Navy disturbs killer whales along with other marine mammals.[187] Killer whales are popular with whale watchers, which may stress the whales and alter their behaviour, particularly if boats approach too closely or block their lines of travel.[188]
+
+The Exxon Valdez oil spill adversely affected killer whales in Prince William Sound and Alaska's Kenai Fjords region. Eleven members (about half) of one resident pod disappeared in the following year. The spill damaged salmon and other prey populations, which in turn damaged local killer whales. By 2009, scientists estimated the AT1 transient population (considered part of a larger population of 346 transients), numbered only seven individuals and had not reproduced since the spill. This population is expected to die out.[189][190]
+
+A 2018 study published in Science found that global killer whale populations are poised to dramatically decline due to exposure to toxic chemical and PCB pollution.[191]
+
+The indigenous peoples of the Pacific Northwest Coast feature killer whales throughout their art, history, spirituality and religion. The Haida regarded killer whales as the most powerful animals in the ocean, and their mythology tells of killer whales living in houses and towns under the sea. According to these myths, they took on human form when submerged, and humans who drowned went to live with them.[192] For the Kwakwaka'wakw, the killer whale was regarded as the ruler of the undersea world, with sea lions for slaves and dolphins for warriors.[192] In Nuu-chah-nulth and Kwakwaka'wakw mythology, killer whales may embody the souls of deceased chiefs.[192] The Tlingit of southeastern Alaska regarded the killer whale as custodian of the sea and a benefactor of humans.[193]
+
+The Maritime Archaic people of Newfoundland also had great respect for killer whales, as evidenced by stone carvings found in a 4,000-year-old burial at the Port au Choix Archaeological Site.[194][195]
+
+In the tales and beliefs of the Siberian Yupik people, killer whales are said to appear as wolves in winter, and wolves as killer whales in summer.[196][197][198][199] Killer whales are believed to assist their hunters in driving walrus.[200] Reverence is expressed in several forms: the boat represents the animal, and a wooden carving hung from the hunter's belt.[198] Small sacrifices such as tobacco or meat are strewn into the sea for them.[200][199]
+
+Indigenous Ainu tribe often referred killer whales in their folklore and myth as Repun Kamuy (God of Sea/Offshore) to bring fortunes (whales) to the coasts, and there had been traditional funerals for stranded or deceased orcas akin to funerals for other animals such as brown bears.[201]
+
+In Western cultures, killer whales were historically feared as dangerous, savage predators.[202] The first written description of a killer whale was given by Pliny the Elder circa AD 70, who wrote, "Orcas (the appearance of which no image can express, other than an enormous mass of savage flesh with teeth) are the enemy of [other kinds of whale]... they charge and pierce them like warships ramming."[203]
+
+Of the very few confirmed attacks on humans by wild killer whales, none have been fatal.[204] In one instance, killer whales tried to tip ice floes on which a dog team and photographer of the Terra Nova Expedition were standing.[205] The sled dogs' barking is speculated to have sounded enough like seal calls to trigger the killer whale's hunting curiosity. In the 1970s, a surfer in California was bitten, and in 2005, a boy in Alaska who was splashing in a region frequented by harbour seals was bumped by a killer whale that apparently misidentified him as prey.[206] Unlike wild killer whales, captive killer whales have made nearly two dozen attacks on humans since the 1970s, some of which have been fatal.[207][208]
+
+Competition with fishermen also led to killer whales being regarded as pests. In the waters of the Pacific Northwest and Iceland, the shooting of killer whales was accepted and even encouraged by governments.[202] As an indication of the intensity of shooting that occurred until fairly recently, about 25% of the killer whales captured in Puget Sound for aquarium through 1970 bore bullet scars.[209] The U.S. Navy claimed to have deliberately killed hundreds of killer whales in Icelandic waters in 1956 with machine guns, rockets, and depth charges.[210][211]
+
+Western attitudes towards killer whales have changed dramatically in recent decades. In the mid-1960s and early 1970s, killer whales came to much greater public and scientific awareness, starting with the first live-capture and display of a killer whale known as Moby Doll, a resident harpooned off Saturna Island in 1964.[202] So little was known at the time, it was nearly two months before the whale's keepers discovered what food (fish) it was willing to eat. To the surprise of those who saw him, Moby Doll was a docile, non-aggressive whale who made no attempts to attack humans.[212]
+
+Between 1964 and 1976, 50 killer whales from the Pacific Northwest were captured for display in aquaria, and public interest in the animals grew. In the 1970s, research pioneered by Michael Bigg led to the discovery of the species' complex social structure, its use of vocal communication, and its extraordinarily stable mother–offspring bonds. Through photo-identification techniques, individuals were named and tracked over decades.[213]
+
+Bigg's techniques also revealed the Pacific Northwest population was in the low hundreds rather than the thousands that had been previously assumed.[202] The southern resident community alone had lost 48 of its members to captivity; by 1976, only 80 remained.[214] In the Pacific Northwest, the species that had unthinkingly been targeted became a cultural icon within a few decades.[181]
+
+The public's growing appreciation also led to growing opposition to whale–keeping in aquarium. Only one whale has been taken in North American waters since 1976. In recent years, the extent of the public's interest in killer whales has manifested itself in several high-profile efforts surrounding individuals. Following the success of the 1993 film Free Willy, the movie's captive star Keiko was returned to the coast of his native Iceland in 2002. The director of the International Marine Mammal Project for the Earth Island Institute, David Phillips, led the efforts to return Keiko to the Iceland waters.[215] Keiko however did not adapt to the harsh climate of the Arctic Ocean, and died a year into his release after contracting pneumonia, at the age of 27.[216] In 2002, the orphan Springer was discovered in Puget Sound, Washington. She became the first whale to be successfully reintegrated into a wild pod after human intervention, crystallizing decades of research into the vocal behaviour and social structure of the region's killer whales.[217] The saving of Springer raised hopes that another young killer whale named Luna, which had become separated from his pod, could be returned to it. However, his case was marked by controversy about whether and how to intervene, and in 2006, Luna was killed by a boat propeller.[218]
+
+The earlier of known records of commercial hunting of killer whales date to the 18th century in Japan. During the 19th and early 20th centuries, the global whaling industry caught immense numbers of baleen and sperm whales, but largely ignored killer whales because of their limited amounts of recoverable oil, their smaller populations, and the difficulty of taking them.[143] Once the stocks of larger species were depleted, killer whales were targeted by commercial whalers in the mid-20th century. Between 1954 and 1997, Japan took 1,178 killer whales (although the Ministry of the Environment claims that there had been domestic catches of about 1,600 whales between late 1940s to 1960s[219]) and Norway took 987.[220] Extensive hunting of killer whales, including an Antarctic catch of 916 in 1979–80 alone, prompted the International Whaling Commission to recommend a ban on commercial hunting of the species pending further research.[220] Today, no country carries out a substantial hunt, although Indonesia and Greenland permit small subsistence hunts (see Aboriginal whaling). Other than commercial hunts, killer whales were hunted along Japanese coasts out of public concern for potential conflicts with fisheries. Such cases include a semi-resident male-female pair in Akashi Strait and Harimanada being killed in the Seto Inland Sea in 1957,[221][222] the killing of five whales from a pod of 11 members that swam into Tokyo Bay in 1970,[223] and a catch record in southern Taiwan in the 1990s.[79][224]
+
+Killer whales have helped humans hunting other whales.[225] One well-known example was the killer whales of Eden, Australia, including the male known as Old Tom. Whalers more often considered them a nuisance, however, as orcas would gather to scavenge meat from the whalers' catch.[225] Some populations, such as in Alaska's Prince William Sound, may have been reduced significantly by whalers shooting them in retaliation.[16]
+
+Whale watching continues to increase in popularity, but may have some problematic impacts on killer whales. Exposure to exhaust gasses from large amounts of vessel traffic are causing concern for the overall health of the 75 remaining southern resident killer whales (SRKWs) left as of early 2019.[226] This population is followed by approximately 20 vessels for 12 hours a day during the months May–September.[227] Researchers discovered that these vessels are in the line of sight for these whales for 98–99.5% of daylight hours.[227] With so many vessels, the air quality around these whales deteriorates and impacts their health. Air pollutants that bind with exhaust fumes are responsible for the activation of the cytochrome P450 1A gene family.[227] Researchers have successfully identified this gene in skin biopsies of live whales and also the lungs of deceased whales. A direct correlation between activation of this gene and the air pollutants can not be made because there are other known factors that will induce the same gene. Vessels can have either wet or dry exhaust systems, with wet exhaust systems leaving more pollutants in the water due to various gas solubility. A modelling study determined that the lowest-observed-adverse-effect-level (LOAEL) of exhaust pollutants was about 12% of the human dose.[227]
+
+As a response to this, in 2017 boats off the British Columbia coast now have a minimum approach distance of 200 metres compared to the previous 100 metres. This new rule complements Washington State's minimum approach zone of 180 metres that has been in effect since 2011. If a whale approaches a vessel it must be placed in neutral until the whale passes. The World Health Organization has set air quality standards in an effort to control the emissions produced by these vessels.[228]
+
+The killer whale's intelligence, trainability, striking appearance, playfulness in captivity and sheer size have made it a popular exhibit at aquaria and aquatic theme parks. From 1976 to 1997, 55 whales were taken from the wild in Iceland, 19 from Japan, and three from Argentina. These figures exclude animals that died during capture. Live captures fell dramatically in the 1990s, and by 1999, about 40% of the 48 animals on display in the world were captive-born.[229]
+
+Organizations such as World Animal Protection and the Whale and Dolphin Conservation campaign against the practice of keeping them in captivity. In captivity, they often develop pathologies, such as the dorsal fin collapse seen in 60–90% of captive males. Captives have vastly reduced life expectancies, on average only living into their 20s.[c] In the wild, females who survive infancy live 46 years on average, and up to 70–80 years in rare cases. Wild males who survive infancy live 31 years on average, and up to 50–60 years.[230] Captivity usually bears little resemblance to wild habitat, and captive whales' social groups are foreign to those found in the wild. Critics claim captive life is stressful due to these factors and the requirement to perform circus tricks that are not part of wild killer whale behaviour, see above.[231] Wild killer whales may travel up to 160 kilometres (100 mi) in a day, and critics say the animals are too big and intelligent to be suitable for captivity.[155] Captives occasionally act aggressively towards themselves, their tankmates, or humans, which critics say is a result of stress.[207] Between 1991 and 2010, the bull orca known as Tilikum was involved in the death of three people, and was featured in the critically acclaimed 2013 film Blackfish.[232] Tilikum lived at SeaWorld from 1992 until his death in 2017.[233][234][235][236][237]
+
+A 2015 study coauthored by staff at SeaWorld and the Minnesota Zoo indicates that there is no significant difference in survivorship between free-ranging and captive killer whales. The authors speculate about the future utility of studying captive populations for the purposes of understanding orca biology and the implications of such research of captive animals in the overall health of both wild and marine park populations.[238]
+
+As of March 2016, SeaWorld has announced that they will be ending their orca breeding program and their theatrical shows. They previously announced, in November 2015, that the shows would be coming to an end in San Diego but it is now to happen in both Orlando and San Antonio as well.[239]
+

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+A sword is a bladed melee weapon intended for cutting or thrusting that is longer than a knife or dagger, consisting of a long blade attached to a hilt. The precise definition of the term varies with the historical epoch or the geographic region under consideration. The blade can be straight or curved. Thrusting swords have a pointed tip on the blade, and tend to be straighter; slashing swords have a sharpened cutting edge on one or both sides of the blade, and are more likely to be curved. Many swords are designed for both thrusting and slashing.
+
+Historically, the sword developed in the Bronze Age, evolving from the dagger; the earliest specimens date to about 1600 BC. The later Iron Age sword remained fairly short and without a crossguard. The spatha, as it developed in the Late Roman army, became the predecessor of the European sword of the Middle Ages, at first adopted as the Migration Period sword, and only in the High Middle Ages, developed into the classical arming sword with crossguard. The word sword continues the Old English, sweord.[1]
+
+The use of a sword is known as swordsmanship or, in a modern context, as fencing. In the Early Modern period, western sword design diverged into roughly two forms, the thrusting swords and the sabers.
+
+Thrusting swords such as the rapier and eventually the smallsword were designed to impale their targets quickly and inflict deep stab wounds. Their long and straight yet light and well balanced design made them highly maneuverable and deadly in a duel but fairly ineffective when used in a slashing or chopping motion. A well aimed lunge and thrust could end a fight in seconds with just the sword's point, leading to the development of a fighting style which closely resembles modern fencing.
+
+The sabre and similar blades such as the cutlass were built more heavily and were more typically used in warfare. Built for slashing and chopping at multiple enemies, often from horseback, the saber's long curved blade and slightly forward weight balance gave it a deadly character all its own on the battlefield. Most sabers also had sharp points and double-edged blades, making them capable of piercing soldier after soldier in a cavalry charge. Sabers continued to see battlefield use until the early 20th century. The US Navy kept tens of thousands of sturdy cutlasses in their armory well into World War II and many were issued to Marines in the Pacific as jungle machetes.
+
+Non-European weapons classified as swords include single-edged weapons such as the Middle Eastern scimitar, the Chinese dao and the related Japanese katana. The Chinese jìan is an example of a non-European double-edged sword, like the European models derived from the double-edged Iron Age sword.
+
+The first weapons that can be described as "swords" date to around 3300 BC. They have been found in Arslantepe, Turkey, are made from arsenical bronze, and are about 60 cm (24 in) long.[2] Some of them are inlaid with silver.
+
+The sword developed from the knife or dagger. A knife is unlike a dagger in that a knife has only one cutting surface, while a dagger has two cutting surfaces. Construction of longer blades became possible during the 3rd millennium BC in the Middle East, first in arsenic copper, then in tin-bronze.
+
+Blades longer than 60 cm (24 in) were rare and not practical until the late Bronze Age because the Young's modulus (stiffness) of bronze is relatively low, and consequently longer blades would bend easily. The development of the sword out of the dagger was gradual; the first weapons that can be classified as swords without any ambiguity are those found in Minoan Crete, dated to about 1700 BC, reaching a total length of more than 100 cm (39 in). These are the "type A" swords of the Aegean Bronze Age.
+
+One of the most important, and longest-lasting, types swords of the European Bronze Age was the Naue II type (named for Julius Naue who first described them), also known as Griffzungenschwert (lit. "grip-tongue sword"). This type first appears in c. the 13th century BC in Northern Italy (or a general Urnfield background), and survives well into the Iron Age, with a life-span of about seven centuries. During its lifetime, metallurgy changed from bronze to iron, but not its basic design.
+
+Naue II swords were exported from Europe to the Aegean, and as far afield as Ugarit, beginning about 1200 BC, i.e. just a few decades before the final collapse of the palace cultures in the Bronze Age collapse.[3] Naue II swords could be as long as 85 cm, but most specimens fall into the 60 to 70 cm range. Robert Drews linked the Naue Type II Swords, which spread from Southern Europe into the Mediterranean, with the Bronze Age collapse.[4] Naue II swords, along with Nordic full-hilted swords, were made with functionality and aesthetics in mind.[5] The hilts of these swords were beautifully crafted and often contained false rivets in order to make the sword more visually appealing. Swords coming from northern Denmark and northern Germany usually contained three or more fake rivets in the hilt.[6]
+
+Sword production in China is attested from the Bronze Age Shang Dynasty.[7] The technology for bronze swords reached its high point during the Warring States period and Qin Dynasty. Amongst the Warring States period swords, some unique technologies were used, such as casting high tin edges over softer, lower tin cores, or the application of diamond shaped patterns on the blade (see sword of Goujian). Also unique for Chinese bronzes is the consistent use of high tin bronze (17–21% tin) which is very hard and breaks if stressed too far, whereas other cultures preferred lower tin bronze (usually 10%), which bends if stressed too far. Although iron swords were made alongside bronze, it was not until the early Han period that iron completely replaced bronze.[8]
+
+In the Indian subcontinent, earliest available Bronze age swords of copper were discovered in the Indus Valley Civilization sites in the northwestern regions of South Asia. Swords have been recovered in archaeological findings throughout the Ganges-Jamuna Doab region of Indian subcontinent, consisting of bronze but more commonly copper.[9] Diverse specimens have been discovered in Fatehgarh, where there are several varieties of hilt.[9] These swords have been variously dated to times between 1700–1400 BC, but were probably used more in the opening centuries of the 1st millennium BC.[9]
+
+Iron became increasingly common from the 13th century BC. Before that the use of swords was less frequent. The iron was not quench-hardened although often containing sufficient carbon, but work-hardened like bronze by hammering. This made them comparable or only slightly better in terms of strength and hardness to bronze swords. They could still bend during use rather than spring back into shape. But the easier production, and the better availability of the raw material for the first time permitted the equipment of entire armies with metal weapons, though Bronze Age Egyptian armies were sometimes fully equipped with bronze weapons.[10]
+
+Ancient swords are often found at burial sites. The sword was often placed on the right side of the corpse. Many times the sword was kept over the corpse. In many late Iron Age graves, the sword and the scabbard were bent at 180 degrees. It was known as killing the sword. Thus they might have considered swords as the most potent and powerful object.[11]
+
+By the time of Classical Antiquity and the Parthian and Sassanid Empires in Iran, iron swords were common. The Greek xiphos and the Roman gladius are typical examples of the type, measuring some 60 to 70 cm (24 to 28 in).[12][13] The late Roman Empire introduced the longer spatha[14] (the term for its wielder, spatharius, became a court rank in Constantinople), and from this time, the term longsword is applied to swords comparatively long for their respective periods.[15]
+
+Swords from the Parthian and Sassanian Empires were quite long, the blades on some late Sassanian swords being just under a metre long.
+
+Swords were also used to administer various physical punishments, such as non-surgical amputation or capital punishment by decapitation. The use of a sword, an honourable weapon, was regarded in Europe since Roman times as a privilege reserved for the nobility and the upper classes.[16]
+
+The Periplus of the Erythraean Sea mentions swords of Indian iron and steel being exported from ancient India to Greece.[17] Blades from the Indian subcontinent made of Damascus steel also found their way into Persia.[17]
+
+In the first millennium BC the Persian armies used a sword that was originally of Scythian design called the akinaka (acinaces). However, the great conquests of the Persians made the sword more famous as a Persian weapon, to the extent that the true nature of the weapon has been lost somewhat as the name Akinaka has been used to refer to whichever form of sword the Persian army favoured at the time.
+
+It is widely believed that the original akinaka was a 35 to 45 cm (14 to 18 inch) double-edged sword. The design was not uniform and in fact identification is made more on the nature of the scabbard than the weapon itself; the scabbard usually has a large, decorative mount allowing it to be suspended from a belt on the wearer's right side. Because of this, it is assumed that the sword was intended to be drawn with the blade pointing downwards ready for surprise stabbing attacks.
+
+In the 12th century, the Seljuq dynasty had introduced the curved shamshir to Persia, and this was in extensive use by the early 16th century.
+
+Chinese iron swords made their first appearance in the later part of the Western Zhou Dynasty, but iron and steel swords were not widely used until the 3rd century BC Han Dynasty.[8] The Chinese Dao (刀 pinyin dāo) is single-edged, sometimes translated as sabre or broadsword, and the Jian (劍 or 剑 pinyin jiàn) is double-edged. The zhanmadao (literally "horse chopping sword"), an extremely long, anti-cavalry sword from the Song dynasty era.
+
+During the Middle Ages sword technology improved, and the sword became a very advanced weapon. The spatha type remained popular throughout the Migration period and well into the Middle Ages. Vendel Age spathas were decorated with Germanic artwork (not unlike the Germanic bracteates fashioned after Roman coins). The Viking Age saw again a more standardized production, but the basic design remained indebted to the spatha.[18]
+
+Around the 10th century, the use of properly quenched hardened and tempered steel started to become much more common than in previous periods. The Frankish 'Ulfberht' blades (the name of the maker inlaid in the blade) were of particularly consistent high quality.[19] Charles the Bald tried to prohibit the export of these swords, as they were used by Vikings in raids against the Franks.
+
+Wootz steel which is also known as Damascus steel was a unique and highly prized steel developed on the Indian subcontinent as early as the 5th century BC. Its properties were unique due to the special smelting and reworking of the steel creating networks of iron carbides described as a globular cementite in a matrix of pearlite. The use of Damascus steel in swords became extremely popular in the 16th and 17th centuries.[nb 1][20]
+
+It was only from the 11th century that Norman swords began to develop the crossguard (quillons). During the Crusades of the 12th to 13th century, this cruciform type of arming sword remained essentially stable, with variations mainly concerning the shape of the pommel. These swords were designed as cutting weapons, although effective points were becoming common to counter improvements in armour, especially the 14th-century change from mail to plate armour.[21]
+
+It was during the 14th century, with the growing use of more advanced armour, that the hand and a half sword, also known as a "bastard sword", came into being. It had an extended grip that meant it could be used with either one or two hands. Though these swords did not provide a full two-hand grip they allowed their wielders to hold a shield or parrying dagger in their off hand, or to use it as a two-handed sword for a more powerful blow.[22]
+
+In the Middle Ages, the sword was often used as a symbol of the word of God. The names given to many swords in mythology, literature, and history reflected the high prestige of the weapon and the wealth of the owner.[23]
+
+From around 1300 to 1500, in concert with improved armour, innovative sword designs evolved more and more rapidly. The main transition was the lengthening of the grip, allowing two-handed use, and a longer blade. By 1400, this type of sword, at the time called langes Schwert (longsword) or spadone, was common, and a number of 15th- and 16th-century Fechtbücher offering instructions on their use survive. Another variant was the specialized armour-piercing swords of the estoc type. The longsword became popular due to its extreme reach and its cutting and thrusting abilities.[24]
+
+The estoc became popular because of its ability to thrust into the gaps between plates of armour.[25] The grip was sometimes wrapped in wire or coarse animal hide to provide a better grip and to make it harder to knock a sword out of the user's hand.[26]
+
+A number of manuscripts covering longsword combat and techniques dating from the 13th–16th centuries exist in German,[27] Italian, and English,[28] providing extensive information on longsword combatives as used throughout this period. Many of these are now readily available online.[27][28]
+
+In the 16th century, the large zweihänder was used by the elite German and Swiss mercenaries known as doppelsöldners.[29] Zweihänder, literally translated, means two-hander. The zweihänder possesses a long blade, as well as a huge guard for protection. It is estimated that some zweihänder swords were over 6 feet (1.8 m) long, with the one ascribed to Frisian warrior Pier Gerlofs Donia being 7 feet (2.13 m) long.[30] The gigantic blade length was perfectly designed for manipulating and pushing away enemy pole-arms, which were major weapons around this time, in both Germany and Eastern Europe. Doppelsöldners also used katzbalgers, which means 'cat-gutter'. The katzbalger's S-shaped guard and 2-foot-long (0.61 m) blade made it perfect for bringing in when the fighting became too close to use a zweihänder.[31]
+
+Civilian use of swords became increasingly common during the late Renaissance, with duels being a preferred way to honourably settle disputes.
+
+The side-sword was a type of war sword used by infantry during the Renaissance of Europe. This sword was a direct descendant of the arming sword. Quite popular between the 16th and 17th centuries, they were ideal for handling the mix of armoured and unarmoured opponents of that time. A new technique of placing one's finger on the ricasso to improve the grip (a practice that would continue in the rapier) led to the production of hilts with a guard for the finger. This sword design eventually led to the development of the civilian rapier, but it was not replaced by it, and the side-sword continued to be used during the rapier's lifetime. As it could be used for both cutting and thrusting, the term cut and thrust sword is sometimes used interchangeably with side-sword.[32] As rapiers became more popular, attempts were made to hybridize the blade, sacrificing the effectiveness found in each unique weapon design. These are still considered side-swords and are sometimes labeled sword rapier or cutting rapier by modern collectors.
+
+Side-swords used in conjunction with bucklers became so popular that it caused the term swashbuckler to be coined. This word stems from the new fighting style of the side-sword and buckler which was filled with much "swashing and making a noise on the buckler".[33]
+
+Within the Ottoman Empire, the use of a curved sabre called the Yatagan started in the mid-16th century. It would become the weapon of choice for many in Turkey and the Balkans.[34]
+
+The sword in this time period was the most personal weapon, the most prestigious, and the most versatile for close combat, but it came to decline in military use as technology, such as the crossbow and firearms changed warfare. However, it maintained a key role in civilian self-defence.[35]
+
+The earliest evidence of curved swords, or scimitars (and other regional variants as the Arabian saif, the Persian shamshir and the Turkic kilij) is from the 9th century, when it was used among soldiers in the Khurasan region of Persia.[36]
+
+Kilij.
+
+Shamshir.
+
+The takoba is a type of broadsword originating in the Sahel, descended from the various Byzantine and Islamic swords used across North Africa. Strongly associated with the Tuaregs, it has a straight double-edged blade measuring about 1 meter in length, usually imported from Europe.[37][38]
+Abyssinian swords related to the Persian shamshir are known as  shotel.[39]
+The Ashanti people adopted swords under the name of akrafena. They are still used today in ceremonies, such as the Odwira festival.[40][41]
+
+As steel technology improved, single-edged weapons became popular throughout Asia. Derived from the Chinese Jian or dao, the Korean hwandudaedo are known from the early medieval Three Kingdoms. Production of the Japanese tachi, a precursor to the katana, is recorded from c. AD 900 (see Japanese sword).[42]
+
+Japan was famous for the swords it forged in the early 13th century for the class of warrior-nobility known as the Samurai. The types of swords used by the Samurai included the ōdachi (extra long field sword), tachi (long cavalry sword), katana (long sword), and wakizashi (shorter companion sword for katana). Japanese swords that pre-date the rise of the samurai caste include the tsurugi (straight double-edged blade) and chokutō (straight one-edged blade).[43] Japanese swordmaking reached the height of its development in the 15th and 16th centuries, when samurai increasingly found a need for a sword to use in closer quarters, leading to the creation of the modern katana.[44]
+
+Western historians have said that Japanese katana were among the finest cutting weapons in world military history.[45][46][47]
+
+In Indonesia, the images of Indian style swords can be found in Hindu gods statues from ancient Java circa 8th to 10th century. However the native types of blade known as kris, parang, klewang and golok were more popular as weapons. These daggers are shorter than sword but longer than common dagger.
+
+In The Philippines, traditional large swords known as the Kampilan and the Panabas were used in combat by the natives. A notable wielder of the kampilan was Lapu-Lapu, the king of Mactan and his warriors who defeated the Spaniards and killed Portuguese explorer Ferdinand Magellan at the Battle of Mactan on 27 April 1521.[48] Traditional swords in the Philippines were immediately banned, but the training in swordsmanship was later hidden from the occupying Spaniards by practices in dances. But because of the banning, Filipinos were forced to use swords that were disguised as farm tools. Bolos and baliswords were used during the revolutions against the colonialists not only because ammunition for guns was scarce, but also for concealability while walking in crowded streets and homes. Bolos were also used by young boys who joined their parents in the revolution and by young girls and their mothers in defending the town while the men were on the battlefields. During the Philippine–American War in events such as the Balangiga Massacre, most of an American company was hacked to death or seriously injured by bolo-wielding guerillas in Balangiga, Samar.[49] When the Japanese took control of the country, several American special operations groups stationed in the Philippines were introduced to the Filipino Martial Arts and swordsmanship, leading to this style reaching America despite the fact that natives were reluctant to allow outsiders in on their fighting secrets.[50]
+
+The Khanda is a double-edge straight sword. It is often featured in religious iconography, theatre and art depicting the ancient history of India. Some communities venerate the weapon as a symbol of Shiva. It is a common weapon in the martial arts in the Indian subcontinent.[51] Khanda often appears in Hindu, Buddhist and Sikh scriptures and art.[52] In Sri Lanka, a unique wind furnace was used to produce the high quality steel. This gave the blade a very hard cutting edge and beautiful patterns. For these reasons it became a very popular trading material.[53]
+
+The Firangi (/fəˈrɪŋɡiː/, derived from the Arabic term for a Western European a "Frank") was a sword type which used blades manufactured in Western Europe and imported by the Portuguese, or made locally in imitation of European blades. Because of its length the firangi is usually regarded as primarily a cavalry weapon. The sword has been especially associated with the Marathas, who were famed for their cavalry. However, the firangi was also widely used by Sikhs and Rajputs.[54]
+
+The Talwar (Hindi: तलवार) is a type of curved sword from India and other countries of the Indian subcontinent, it was adopted by communities such as Rajputs, Sikhs and Marathas, who favored the sword as their main weapon. It became more widespread in the medieval era.[55][56]
+
+The Urumi (Tamil: சுருள் பட்டாக்கத்தி surul pattai, lit. curling blade; Sinhala: එතුණු කඩුව ethunu kaduwa; Hindi: aara) is a "sword" with a flexible whip-like blade.[57]
+
+Talwar
+
+Pata
+
+Firangi
+
+A single-edged type of sidearm used by the Hussites was popularized in 16th-century Germany under its Czech name Dusack, also known as Säbel auf Teutsch gefasst ("sabre fitted in the German manner").[58] A closely related weapon is the schnepf or Swiss sabre used in Early Modern Switzerland.[59]
+
+The cut-and-thrust mortuary sword was used after 1625 by cavalry during the English Civil War. This (usually) two-edged sword sported a half-basket hilt with a straight blade some 90–105 cm long. Later in the 17th century, the swords used by cavalry became predominantly single-edged. The so-called walloon sword (épée wallone)[60] was common in the Thirty Years' War and Baroque era.[61] Its hilt was ambidextrous with shell-guards and knuckle-bow that inspired 18th century continental hunting hangers.[62] Following their campaign in the Netherlands in 1672, the French began producing this weapon as their first regulation sword.[63] Weapons of this design were also issued to the Swedish army from the time of Gustavus Adolphus until as late as the 1850s.[64]
+
+The rapier is believed to have evolved either from the Spanish espada ropera or from the swords of the Italian nobility somewhere in the later part of the 16th century.[65][66] The rapier differed from most earlier swords in that it was not a military weapon but a primarily civilian sword. Both the rapier and the Italian schiavona developed the crossguard into a basket-shaped guard for hand protection.[67]
+During the 17th and 18th centuries, the shorter smallsword became an essential fashion accessory in European countries and the New World, though in some places such as the Scottish Highlands large swords as the basket-hilted broadsword were preferred, and most wealthy men and military officers carried one slung from a belt. Both the smallsword and the rapier remained popular dueling swords well into the 18th century.[68]
+
+As the wearing of swords fell out of fashion, canes took their place in a gentleman's wardrobe. This developed to the gentlemen in the Victorian era to use the umbrella. Some examples of canes—those known as sword canes or swordsticks—incorporate a concealed blade. The French martial art la canne developed to fight with canes and swordsticks and has now evolved into a sport. The English martial art singlestick is very similar.
+With the rise of the pistol duel, the duelling sword fell out of fashion long before the practice of duelling itself. By about 1770, English duelists enthusiastically adopted the pistol, and sword duels dwindled.[69] However, the custom of duelling with epées persisted well into the 20th century in France. Such modern duels were not fought to the death; the duellists' aim was instead merely to draw blood from the opponent's sword arm.[70]
+
+Towards the end of its useful life, the sword served more as a weapon of self-defence than for use on the battlefield, and the military importance of swords steadily decreased during the Modern Age. Even as a personal sidearm, the sword began to lose its preeminence in the early 19th century, reflecting the development of reliable handguns.[35]
+
+However, swords were still normally carried in combat by cavalrymen and by officers of other branches throughout the 19th and early 20th centuries, both in colonial and European warfare. For example, during the Aceh War the Acehnese Klewangs, a sword similar to the machete, proved very effective in close quarters combat with Dutch troops, leading the Royal Netherlands East Indies Army to adopt a heavy cutlass, also called klewang (very similar in appearance to the US Navy Model 1917 Cutlass) to counter it. Mobile troops armed with carbines and klewangs succeeded in suppressing Aceh resistance where traditional infantry with rifle and bayonet had failed. From that time on until the 1950s the Royal Dutch East Indies Army, Royal Dutch Army, Royal Dutch Navy and Dutch police used these cutlasses called Klewang.[71][72]
+
+Swords continued in general peacetime use by cavalry of most armies during the years prior to World War I. The British Army formally adopted a completely new design of cavalry sword in 1908, almost the last change in British Army weapons before the outbreak of the war.[73] At the outbreak of World War I infantry officers in all combatant armies then involved (French, German, British, Austro-Hungarian, Russian, Belgian and Serbian) still carried swords as part of their field equipment. On mobilization in August 1914 all serving British Army officers were required to have their swords sharpened as the only peacetime use of the weapon had been for saluting on parade.[74] The high visibility and limited practical use of the sword however led to it being abandoned within weeks, although most cavalry continued to carry sabres throughout the war. While retained as a symbol of rank and status by at least senior officers of infantry, artillery and other branches the sword was usually left with non-essential bagage when units reached the front line.[75] It was not until the late 1920s and early 1930s that this historic weapon was finally discarded for all but ceremonial purposes by most remaining horse mounted regiments of Europe and the Americas.
+
+In China troops used the long anti-cavalry Miao dao well into the Second Sino-Japanese War. The last units of British heavy cavalry switched to using armoured vehicles as late as 1938. Swords and other dedicated melee weapons were used occasionally by many countries during World War II, but typically as a secondary weapon as they were outclassed by coexisting firearms.[76][77][78] A notable exception was the Imperial Japanese Army where, for cultural reasons, all officers and warrant officers carried the Type 94 shin-gunto ("new miltary sword") into battle from 1934 until 1945.[79]
+
+Swords are commonly worn as a ceremonial item by officers in many military and naval services throughout the world. Occasions to wear swords include any event in dress uniforms where the rank-and-file carry arms: parades, reviews, courts-martial, tattoos, and changes of command. They are also commonly worn for officers' weddings, and when wearing dress uniforms to church—although they are rarely actually worn in the church itself.
+
+In the British forces they are also worn for any appearance at Court. In the United States, every Naval officer at or above the rank of Lieutenant Commander is required to own a sword, which can be prescribed for any formal outdoor ceremonial occasion; they are normally worn for changes of command and parades. For some Navy parades, cutlasses are issued to Petty Officers and Chief Petty Officers.
+
+In the U.S. Marine Corps every officer must own a sword, which is prescribed for formal parades and other ceremonies where dress uniforms are worn and the rank-and-file are under arms. On these occasions depending on their billet, Marine Non-Commissioned Officers (E-6 and above) may also be required to carry swords, which have hilts of a pattern similar to U.S. Naval officers' swords but are actually sabres. The USMC Model 1859 NCO Sword is the longest continuously-issued edged weapon in the U.S. inventory
+
+The Marine officer swords are of the Mameluke pattern which was adopted in 1825 in recognition of the Marines' key role in the capture of the Tripolitan city of Derna during the First Barbary War.[80] Taken out of issue for approximately 20 years from 1855 until 1875, it was restored to service in the year of the Corps' centennial and has remained in issue since.
+
+The production of replicas of historical swords originates with 19th-century historicism.[81] Contemporary replicas can range from cheap factory produced look-alikes to exact recreations of individual artifacts, including an approximation of the historical production methods.
+
+Some kinds of swords are still commonly used today as weapons, often as a side arm for military infantry. The Japanese katana, wakizashi and tanto are carried by some infantry and officers in Japan and other parts of Asia and the kukri is the official melee weapon for Nepal. Other swords in use today are the sabre, the scimitar, the shortsword and the machete.[82]
+
+The sword consists of the blade and the hilt.
+The term scabbard applies to the cover for the sword blade when not in use.
+
+There is considerable variation in the detailed design of sword blades. The diagram opposite shows a typical Medieval European sword.
+
+Early iron blades have rounded points due to the limited metallurgy of the time. These were still effective for thrusting against lightly armoured opponents. As armour advanced, blades were made narrower, stiffer and sharply pointed to defeat the armour by thrusting.
+
+Dedicated cutting blades are wide and thin, and often have grooves known as fullers which lighten the blade at the cost of some of the blade's stiffness. The edges of a cutting sword are almost parallel. Blades oriented for the thrust have thicker blades, sometimes with a distinct midrib for increased stiffness, with a strong taper and an acute point.  The geometry of a cutting sword blade allows for acute edge angles. An edge with an acuter angle is more inclined to degrade quickly in combat situations than an edge with a more obtuse angle. Also, an acute edge angle is not the primary factor of a blade's sharpness.[85]
+
+The part of the blade between the center of percussion (CoP) and the point is called the foible (weak) of the blade, and that between the center of balance (CoB) and the hilt is the forte (strong). The section in between the CoP and the CoB is the middle.
+
+The ricasso or shoulder identifies a short section of blade immediately below the guard that is left completely unsharpened. Many swords have no ricasso. On some large weapons, such as the German Zweihänder, a metal cover surrounded the ricasso, and a swordsman might grip it in one hand to wield the weapon more easily in close-quarter combat.[31]
+The ricasso normally bears the maker's mark.
+
+The tang is the extension of the blade to which the hilt is fitted.
+
+On Japanese blades, the maker's mark appears on the tang under the grip.[86]
+
+The hilt is the collective term for the parts allowing for the handling and control of the blade; these consist of the grip, the pommel, and a simple or elaborate guard, which in post-Viking Age swords could consist of only a crossguard (called a cruciform hilt or quillons). The pommel was originally designed as a stop to prevent the sword slipping from the hand. From around the 11th century onward it became a counterbalance to the blade, allowing a more fluid style of fighting.[dubious  – discuss][87]
+It can also be used as a blunt instrument at close range, and its weight affects the centre of percussion. In later times a sword knot or tassel was sometimes added. By the 17th century, with the growing use of firearms and the accompanying decline in the use of armour, many rapiers and dueling swords had developed elaborate basket hilts, which protect the palm of the wielder and rendered the gauntlet obsolete.[88]
+
+In late medieval and Renaissance era European swords, a flap of leather called the chappe or rain guard was attached to a sword's crossguard at the base of the hilt to protect the mouth of the scabbard and prevent water from entering.[89]
+
+Common accessories to the sword include the scabbard, as well as the 'sword belt'.
+
+Sword typology is based on morphological criteria on one hand (blade shape (cross-section, taper, and length), shape and size of the hilt and pommel)
+and age and place of origin on the other (Bronze Age, Iron Age, European (medieval, early modern, modern), Asian).
+
+The relatively comprehensive Oakeshott typology was created by historian and illustrator Ewart Oakeshott as a way to define and catalogue European swords of the medieval period based on physical form, including blade shape and hilt configuration. The typology also focuses on the smaller, and in some cases contemporary, single-handed swords such as the arming sword.[67]
+
+As noted above, the terms longsword, broad sword, great sword, and Gaelic claymore are used relative to the era under consideration, and each term designates a particular type of sword.
+
+In most Asian countries, a sword (jian 劍, geom (검), ken/tsurugi (剣), pedang) is a double-edged straight-bladed weapon, while a knife or saber (dāo 刀, do (도), to/katana (刀), pisau, golok) refers to a single-edged object.
+
+In Sikh history, the sword is held in very high esteem. A single-edged sword is called a kirpan, and its double-edged counterpart a khanda or tega.[92]
+
+The South Indian churika is a handheld double-edged sword traditionally used in the Malabar region of Kerala. It is also worshipped as the weapon of Vettakkorumakan, the hunter god in Hinduism.
+
+European terminology does give generic names for single-edged and double-edged blades but refers to specific types with the term 'sword' covering them all. For example, the backsword may be so called because it is single-edged but the falchion which is also single-edged is given its own specific name.[93]
+
+A two-handed sword is any sword that usually requires two hands to wield, or more specifically the very large swords of the 16th century.[87]
+
+Throughout history two-handed swords have generally been less common than their one-handed counterparts, one exception being their common use in Japan.
+
+A Hand and a half sword, colloquially known as a "bastard sword", was a sword with an extended grip and sometimes pommel so that it could be used with either one or two hands. Although these swords may not provide a full two-hand grip, they allowed its wielders to hold a shield or parrying dagger in their off hand, or to use it as a two-handed sword for a more powerful blow.[26] These should not be confused with a longsword, two-handed sword, or Zweihänder, which were always intended to be used with two hands.
+
+In fantasy, magic swords often appear, based on their use in myth and legend. The science fiction counterpart to these is known as an energy sword (sometimes also referred to as a "beam sword" or "laser sword"), a sword whose blade consists of, or is augmented by, concentrated energy. A well known example of this type of sword is the lightsaber, shown in the Star Wars franchise.
+
+Media related to Swords at Wikimedia Commons

en/1785.html.txt

@@ -0,0 +1,67 @@
+
+
+A spice is a seed, fruit, root, bark, or other plant substance primarily used for flavoring or coloring food. Spices are distinguished from herbs, which are the leaves, flowers, or stems of plants used for flavoring or as a garnish. Spices are sometimes used in medicine, religious rituals, cosmetics or perfume production.[example  needed]
+
+The spice trade developed throughout the Indian subcontinent[1] and Middle East by at earliest 2000 BCE with cinnamon and black pepper, and in East Asia with herbs and pepper. The Egyptians used herbs for mummification and their demand for exotic spices and herbs helped stimulate world trade. The word spice comes from the Old French word espice, which became epice, and which came from the Latin root spec, the noun referring to "appearance, sort, kind": species has the same root. By 1000 BCE, medical systems based upon herbs could be found in China, Korea, and India. Early uses were connected with magic, medicine, religion, tradition, and preservation.[2]
+
+Cloves were used in Mesopotamia by 1700 BCE.[note 1] The ancient Indian epic Ramayana mentions cloves. The Romans had cloves in the 1st century CE, as Pliny the Elder wrote about them.[4]
+
+The earliest written records of spices come from ancient Egyptian, Chinese, and Indian cultures. The Ebers Papyrus from Early Egyptians that dates from 1550 B.C.E. describes some eight hundred different medicinal remedies and numerous medicinal procedures.[5]
+
+Historians believe that nutmeg, which originates from the Banda Islands in Southeast Asia, was introduced to Europe in the 6th century BCE.[6]
+
+Indonesian merchants traveled around China, India, the Middle East, and the east coast of Africa. Arab merchants facilitated the routes through the Middle East and India. This resulted in the Egyptian port city of Alexandria being the main trading center for spices. The most important discovery prior to the European spice trade were the monsoon winds (40 CE). Sailing from Eastern spice cultivators to Western European consumers gradually replaced the land-locked spice routes once facilitated by the Middle East Arab caravans.[2]
+
+In the story of Genesis, Joseph was sold into slavery by his brothers to spice merchants. In the biblical poem Song of Solomon, the male speaker compares his beloved to many forms of spices.
+
+Spices were among the most demanded and expensive products available in Europe in the Middle Ages,[5] the most common being black pepper, cinnamon (and the cheaper alternative cassia), cumin, nutmeg, ginger and cloves. Given medieval medicine's main theory of humorism, spices and herbs were indispensable to balance "humors" in food,[6] a daily basis for good health at a time of recurrent pandemics. In addition to being desired by those using medieval medicine, the European elite also craved spices in the Middle Ages. An example of the European aristocracy's demand for spice comes from the King of Aragon, who invested substantial resources into bringing back spices to Spain in the 12th century. He was specifically looking for spices to put in wine, and was not alone among European monarchs at the time to have such a desire for spice.[7]
+
+Spices were all imported from plantations in Asia and Africa, which made them expensive. From the 8th until the 15th century, the Republic of Venice had the monopoly on spice trade with the Middle East, and along with it the neighboring Italian maritime republics and city-states. The trade made the region rich. It has been estimated that around 1,000 tons of pepper and 1,000 tons of the other common spices were imported into Western Europe each year during the Late Middle Ages. The value of these goods was the equivalent of a yearly supply of grain for 1.5 million people.[8] The most exclusive was saffron, used as much for its vivid yellow-red color as for its flavor. Spices that have now fallen into obscurity in European cuisine include grains of paradise, a relative of cardamom which mostly replaced pepper in late medieval north French cooking, long pepper, mace, spikenard, galangal and cubeb.
+
+Spain and Portugal were interested in seeking new routes to trade in spices and other valuable products from Asia. The control of trade routes and the spice-producing regions were the main reasons that Portuguese navigator Vasco da Gama sailed to India in 1499.[8] When da Gama discovered the pepper market in India, he was able to secure peppers for a much cheaper price than the ones demanded by Venice.[7] At around the same time, Christopher Columbus returned from the New World.  He described to investors new spices available there.[citation needed]
+
+Another source of competition in the spice trade during the 15th and 16th century was the Ragusans from the maritime republic of Dubrovnik in southern Croatia.[9]
+
+The military prowess of Afonso de Albuquerque (1453–1515) allowed the Portuguese to take control of the sea routes to India. In 1506, he took the island of Socotra in the mouth of the Red Sea and, in 1507, Ormuz in the Persian Gulf. Since becoming the viceroy of the Indies, he took Goa in India in 1510, and Malacca on the Malay peninsula in 1511. The Portuguese could now trade directly with Siam, China, and the Maluku Islands.
+
+With the discovery of the New World came new spices, including allspice, chili peppers, vanilla, and chocolate. This development kept the spice trade, with America as a late comer with its new seasonings, profitable well into the 19th century.[citation needed]
+
+Spices are primarily used as food flavoring. They are also used to perfume cosmetics and incense[10]. At various periods, many spices have been believed to have medicinal value. Finally, since they are expensive, rare, and exotic commodities, their conspicuous consumption has often been a symbol of wealth and social class.[11]
+
+It is often claimed that spices were used either as food preservatives or to mask the taste of spoiled meat, especially in the Middle Ages.[12] This is false.[13][14][15] In fact, spices are rather ineffective as preservatives as compared to salting, smoking, pickling, or drying, and are ineffective in covering the taste of spoiled meat.[11] Moreover, spices have always been comparatively expensive: in 15th century Oxford, a whole pig cost about the same as a pound of the cheapest spice, pepper.[11] There is also no evidence of such use from contemporary cookbooks: "Old cookbooks make it clear that spices weren't used as a preservative. They typically suggest adding spices toward the end of the cooking process, where they could have no preservative effect whatsoever."[16] In fact, Cristoforo di Messisbugo suggested in the 16th century that pepper may speed up spoilage.[16]
+
+Though some spices have antimicrobial properties in vitro,[17] pepper—by far the most common spice—is relatively ineffective, and in any case, salt, which is far cheaper, is also far more effective.[16]
+
+A spice may be available in several forms: fresh, whole dried, or pre-ground dried. Generally, spices are dried. Spices may be ground into a powder for convenience. A whole dried spice has the longest shelf life, so it can be purchased and stored in larger amounts, making it cheaper on a per-serving basis. A fresh spice, such as ginger, is usually more flavorful than its dried form, but fresh spices are more expensive and have a much shorter shelf life. Some spices are not always available either fresh or whole, for example turmeric, and often must be purchased in ground form. Small seeds, such as fennel and mustard seeds, are often used both whole and in powder form.
+
+To grind a whole spice, the classic tool is mortar and pestle. Less labor-intensive tools are more common now:  a microplane or fine grater can be used to grind small amounts; a coffee grinder[note 2] is useful for larger amounts. A frequently used spice such as black pepper may merit storage in its own hand grinder or mill.
+
+The flavor of a spice is derived in part from compounds (volatile oils) that oxidize or evaporate when exposed to air. Grinding a spice greatly increases its surface area and so increases the rates of oxidation and evaporation. Thus, flavor is maximized by storing a spice whole and grinding when needed. The shelf life of a whole dry spice is roughly two years; of a ground spice roughly six months.[18] The "flavor life" of a ground spice can be much shorter.[note 3] Ground spices are better stored away from light.[note 4]
+
+Some flavor elements in spices are soluble in water; many are soluble in oil or fat. As a general rule, the flavors from a spice take time to infuse into the food so spices are added early in preparation. This contrasts to herbs which are usually added late in preparation.[18]
+
+A study by the Food and Drug Administration of shipments of spices to the United States during fiscal years 2007-2009 showed about 7% of the shipments were contaminated by Salmonella bacteria, some of it antibiotic-resistant.[19] As most spices are cooked before being served salmonella contamination often has no effect, but some spices, particularly pepper, are often eaten raw and present at table for convenient use. Shipments from Mexico and India, a major producer, were the most frequently contaminated.[20] However, with newly developed radiation sterilization methods, the risk of Salmonella contamination is now lower.[21]
+
+Because they tend to have strong flavors and are used in small quantities, spices tend to add few calories to food, even though many spices, especially those made from seeds, contain high portions of fat, protein, and carbohydrate by weight. However, when used in larger quantity, spices can also contribute a substantial amount of minerals and other micronutrients, including iron, magnesium, calcium, and many others, to the diet. For example, a teaspoon of paprika contains about 1133 IU of Vitamin A, which is over 20% of the recommended daily allowance specified by the US FDA.[22]
+
+Most herbs and spices have substantial antioxidant activity, owing primarily to phenolic compounds, especially flavonoids, which influence nutrition through many pathways, including affecting the absorption of other nutrients. One study found cumin and fresh ginger to be highest in antioxidant activity.[23]
+
+India contributes 75% of global spice production.
+
+The International Organization for Standardization addresses spices and condiments, along with related food additives, as part of the International Classification for Standards 67.220 series.[25]
+
+The Indian Institute of Spices Research in Kozhikode, Kerala, is devoted exclusively to conducting research for ten spice crops: black pepper, cardamom, cinnamon, clove, garcinia, ginger, nutmeg, paprika, turmeric, and vanilla.
+
+The Gato Negro café and spice shop (Buenos Aires, Argentina)
+
+A spice shop selling a variety of spices in Iran
+
+Night spice shop in Casablanca, Morocco.
+
+A spice shop in Taliparamba, India
+
+Spices sold in Taliparamba, India
+
+Books
+
+Articles

en/1786.html.txt

@@ -0,0 +1,67 @@
+
+
+A spice is a seed, fruit, root, bark, or other plant substance primarily used for flavoring or coloring food. Spices are distinguished from herbs, which are the leaves, flowers, or stems of plants used for flavoring or as a garnish. Spices are sometimes used in medicine, religious rituals, cosmetics or perfume production.[example  needed]
+
+The spice trade developed throughout the Indian subcontinent[1] and Middle East by at earliest 2000 BCE with cinnamon and black pepper, and in East Asia with herbs and pepper. The Egyptians used herbs for mummification and their demand for exotic spices and herbs helped stimulate world trade. The word spice comes from the Old French word espice, which became epice, and which came from the Latin root spec, the noun referring to "appearance, sort, kind": species has the same root. By 1000 BCE, medical systems based upon herbs could be found in China, Korea, and India. Early uses were connected with magic, medicine, religion, tradition, and preservation.[2]
+
+Cloves were used in Mesopotamia by 1700 BCE.[note 1] The ancient Indian epic Ramayana mentions cloves. The Romans had cloves in the 1st century CE, as Pliny the Elder wrote about them.[4]
+
+The earliest written records of spices come from ancient Egyptian, Chinese, and Indian cultures. The Ebers Papyrus from Early Egyptians that dates from 1550 B.C.E. describes some eight hundred different medicinal remedies and numerous medicinal procedures.[5]
+
+Historians believe that nutmeg, which originates from the Banda Islands in Southeast Asia, was introduced to Europe in the 6th century BCE.[6]
+
+Indonesian merchants traveled around China, India, the Middle East, and the east coast of Africa. Arab merchants facilitated the routes through the Middle East and India. This resulted in the Egyptian port city of Alexandria being the main trading center for spices. The most important discovery prior to the European spice trade were the monsoon winds (40 CE). Sailing from Eastern spice cultivators to Western European consumers gradually replaced the land-locked spice routes once facilitated by the Middle East Arab caravans.[2]
+
+In the story of Genesis, Joseph was sold into slavery by his brothers to spice merchants. In the biblical poem Song of Solomon, the male speaker compares his beloved to many forms of spices.
+
+Spices were among the most demanded and expensive products available in Europe in the Middle Ages,[5] the most common being black pepper, cinnamon (and the cheaper alternative cassia), cumin, nutmeg, ginger and cloves. Given medieval medicine's main theory of humorism, spices and herbs were indispensable to balance "humors" in food,[6] a daily basis for good health at a time of recurrent pandemics. In addition to being desired by those using medieval medicine, the European elite also craved spices in the Middle Ages. An example of the European aristocracy's demand for spice comes from the King of Aragon, who invested substantial resources into bringing back spices to Spain in the 12th century. He was specifically looking for spices to put in wine, and was not alone among European monarchs at the time to have such a desire for spice.[7]
+
+Spices were all imported from plantations in Asia and Africa, which made them expensive. From the 8th until the 15th century, the Republic of Venice had the monopoly on spice trade with the Middle East, and along with it the neighboring Italian maritime republics and city-states. The trade made the region rich. It has been estimated that around 1,000 tons of pepper and 1,000 tons of the other common spices were imported into Western Europe each year during the Late Middle Ages. The value of these goods was the equivalent of a yearly supply of grain for 1.5 million people.[8] The most exclusive was saffron, used as much for its vivid yellow-red color as for its flavor. Spices that have now fallen into obscurity in European cuisine include grains of paradise, a relative of cardamom which mostly replaced pepper in late medieval north French cooking, long pepper, mace, spikenard, galangal and cubeb.
+
+Spain and Portugal were interested in seeking new routes to trade in spices and other valuable products from Asia. The control of trade routes and the spice-producing regions were the main reasons that Portuguese navigator Vasco da Gama sailed to India in 1499.[8] When da Gama discovered the pepper market in India, he was able to secure peppers for a much cheaper price than the ones demanded by Venice.[7] At around the same time, Christopher Columbus returned from the New World.  He described to investors new spices available there.[citation needed]
+
+Another source of competition in the spice trade during the 15th and 16th century was the Ragusans from the maritime republic of Dubrovnik in southern Croatia.[9]
+
+The military prowess of Afonso de Albuquerque (1453–1515) allowed the Portuguese to take control of the sea routes to India. In 1506, he took the island of Socotra in the mouth of the Red Sea and, in 1507, Ormuz in the Persian Gulf. Since becoming the viceroy of the Indies, he took Goa in India in 1510, and Malacca on the Malay peninsula in 1511. The Portuguese could now trade directly with Siam, China, and the Maluku Islands.
+
+With the discovery of the New World came new spices, including allspice, chili peppers, vanilla, and chocolate. This development kept the spice trade, with America as a late comer with its new seasonings, profitable well into the 19th century.[citation needed]
+
+Spices are primarily used as food flavoring. They are also used to perfume cosmetics and incense[10]. At various periods, many spices have been believed to have medicinal value. Finally, since they are expensive, rare, and exotic commodities, their conspicuous consumption has often been a symbol of wealth and social class.[11]
+
+It is often claimed that spices were used either as food preservatives or to mask the taste of spoiled meat, especially in the Middle Ages.[12] This is false.[13][14][15] In fact, spices are rather ineffective as preservatives as compared to salting, smoking, pickling, or drying, and are ineffective in covering the taste of spoiled meat.[11] Moreover, spices have always been comparatively expensive: in 15th century Oxford, a whole pig cost about the same as a pound of the cheapest spice, pepper.[11] There is also no evidence of such use from contemporary cookbooks: "Old cookbooks make it clear that spices weren't used as a preservative. They typically suggest adding spices toward the end of the cooking process, where they could have no preservative effect whatsoever."[16] In fact, Cristoforo di Messisbugo suggested in the 16th century that pepper may speed up spoilage.[16]
+
+Though some spices have antimicrobial properties in vitro,[17] pepper—by far the most common spice—is relatively ineffective, and in any case, salt, which is far cheaper, is also far more effective.[16]
+
+A spice may be available in several forms: fresh, whole dried, or pre-ground dried. Generally, spices are dried. Spices may be ground into a powder for convenience. A whole dried spice has the longest shelf life, so it can be purchased and stored in larger amounts, making it cheaper on a per-serving basis. A fresh spice, such as ginger, is usually more flavorful than its dried form, but fresh spices are more expensive and have a much shorter shelf life. Some spices are not always available either fresh or whole, for example turmeric, and often must be purchased in ground form. Small seeds, such as fennel and mustard seeds, are often used both whole and in powder form.
+
+To grind a whole spice, the classic tool is mortar and pestle. Less labor-intensive tools are more common now:  a microplane or fine grater can be used to grind small amounts; a coffee grinder[note 2] is useful for larger amounts. A frequently used spice such as black pepper may merit storage in its own hand grinder or mill.
+
+The flavor of a spice is derived in part from compounds (volatile oils) that oxidize or evaporate when exposed to air. Grinding a spice greatly increases its surface area and so increases the rates of oxidation and evaporation. Thus, flavor is maximized by storing a spice whole and grinding when needed. The shelf life of a whole dry spice is roughly two years; of a ground spice roughly six months.[18] The "flavor life" of a ground spice can be much shorter.[note 3] Ground spices are better stored away from light.[note 4]
+
+Some flavor elements in spices are soluble in water; many are soluble in oil or fat. As a general rule, the flavors from a spice take time to infuse into the food so spices are added early in preparation. This contrasts to herbs which are usually added late in preparation.[18]
+
+A study by the Food and Drug Administration of shipments of spices to the United States during fiscal years 2007-2009 showed about 7% of the shipments were contaminated by Salmonella bacteria, some of it antibiotic-resistant.[19] As most spices are cooked before being served salmonella contamination often has no effect, but some spices, particularly pepper, are often eaten raw and present at table for convenient use. Shipments from Mexico and India, a major producer, were the most frequently contaminated.[20] However, with newly developed radiation sterilization methods, the risk of Salmonella contamination is now lower.[21]
+
+Because they tend to have strong flavors and are used in small quantities, spices tend to add few calories to food, even though many spices, especially those made from seeds, contain high portions of fat, protein, and carbohydrate by weight. However, when used in larger quantity, spices can also contribute a substantial amount of minerals and other micronutrients, including iron, magnesium, calcium, and many others, to the diet. For example, a teaspoon of paprika contains about 1133 IU of Vitamin A, which is over 20% of the recommended daily allowance specified by the US FDA.[22]
+
+Most herbs and spices have substantial antioxidant activity, owing primarily to phenolic compounds, especially flavonoids, which influence nutrition through many pathways, including affecting the absorption of other nutrients. One study found cumin and fresh ginger to be highest in antioxidant activity.[23]
+
+India contributes 75% of global spice production.
+
+The International Organization for Standardization addresses spices and condiments, along with related food additives, as part of the International Classification for Standards 67.220 series.[25]
+
+The Indian Institute of Spices Research in Kozhikode, Kerala, is devoted exclusively to conducting research for ten spice crops: black pepper, cardamom, cinnamon, clove, garcinia, ginger, nutmeg, paprika, turmeric, and vanilla.
+
+The Gato Negro café and spice shop (Buenos Aires, Argentina)
+
+A spice shop selling a variety of spices in Iran
+
+Night spice shop in Casablanca, Morocco.
+
+A spice shop in Taliparamba, India
+
+Spices sold in Taliparamba, India
+
+Books
+
+Articles

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+An earthquake (also known as a quake, tremor or temblor) is the shaking of the surface of the Earth resulting from a sudden release of energy in the Earth's lithosphere that creates seismic waves. Earthquakes can range in size from those that are so weak that they cannot be felt to those violent enough to propel objects and people into the air, and wreak destruction across entire cities. The seismicity, or seismic activity, of an area is the frequency, type, and size of earthquakes experienced over a period of time. The word tremor is also used for non-earthquake seismic rumbling.
+
+At the Earth's surface, earthquakes manifest themselves by shaking and displacing or disrupting the ground. When the epicenter of a large earthquake is located offshore, the seabed may be displaced sufficiently to cause a tsunami. Earthquakes can also trigger landslides and occasionally, volcanic activity.
+
+In its most general sense, the word earthquake is used to describe any seismic event—whether natural or caused by humans—that generates seismic waves. Earthquakes are caused mostly by rupture of geological faults but also by other events such as volcanic activity, landslides, mine blasts, and nuclear tests. An earthquake's point of initial rupture is called its hypocenter or focus. The epicenter is the point at ground level directly above the hypocenter.
+
+Tectonic earthquakes occur anywhere in the earth where there is sufficient stored elastic strain energy to drive fracture propagation along a fault plane. The sides of a fault move past each other smoothly and aseismically only if there are no irregularities or asperities along the fault surface that increase the frictional resistance. Most fault surfaces do have such asperities, which leads to a form of stick-slip behavior. Once the fault has locked, continued relative motion between the plates leads to increasing stress and therefore, stored strain energy in the volume around the fault surface. This continues until the stress has risen sufficiently to break through the asperity, suddenly allowing sliding over the locked portion of the fault, releasing the stored energy.[1] This energy is released as a combination of radiated elastic strain seismic waves,[2] frictional heating of the fault surface, and cracking of the rock, thus causing an earthquake. This process of gradual build-up of strain and stress punctuated by occasional sudden earthquake failure is referred to as the elastic-rebound theory. It is estimated that only 10 percent or less of an earthquake's total energy is radiated as seismic energy. Most of the earthquake's energy is used to power the earthquake fracture growth or is converted into heat generated by friction. Therefore, earthquakes lower the Earth's available elastic potential energy and raise its temperature, though these changes are negligible compared to the conductive and convective flow of heat out from the Earth's deep interior.[3]
+
+There are three main types of fault, all of which may cause an interplate earthquake: normal, reverse (thrust), and strike-slip. Normal and reverse faulting are examples of dip-slip, where the displacement along the fault is in the direction of dip and where movement on them involves a vertical component. Normal faults occur mainly in areas where the crust is being extended such as a divergent boundary. Reverse faults occur in areas where the crust is being shortened such as at a convergent boundary. Strike-slip faults are steep structures where the two sides of the fault slip horizontally past each other; transform boundaries are a particular type of strike-slip fault. Many earthquakes are caused by movement on faults that have components of both dip-slip and strike-slip; this is known as oblique slip.
+
+Reverse faults, particularly those along convergent plate boundaries, are associated with the most powerful earthquakes, megathrust earthquakes, including almost all of those of magnitude 8 or more. Strike-slip faults, particularly continental transforms, can produce major earthquakes up to about magnitude 8. Earthquakes associated with normal faults are generally less than magnitude 7. For every unit increase in magnitude, there is a roughly thirtyfold increase in the energy released. For instance, an earthquake of magnitude 6.0 releases approximately 32 times more energy than a 5.0 magnitude earthquake and a 7.0 magnitude earthquake releases 1,000 times more energy than a 5.0 magnitude of earthquake. An 8.6 magnitude earthquake releases the same amount of energy as 10,000 atomic bombs like those used in World War II.[4]
+
+This is so because the energy released in an earthquake, and thus its magnitude, is proportional to the area of the fault that ruptures[5] and the stress drop. Therefore, the longer the length and the wider the width of the faulted area, the larger the resulting magnitude. The topmost, brittle part of the Earth's crust, and the cool slabs of the tectonic plates that are descending down into the hot mantle, are the only parts of our planet that can store elastic energy and release it in fault ruptures. Rocks hotter than about 300 °C (572 °F) flow in response to stress; they do not rupture in earthquakes.[6][7] The maximum observed lengths of ruptures and mapped faults (which may break in a single rupture) are approximately 1,000 km (620 mi). Examples are the earthquakes in Alaska (1957), Chile (1960), and Sumatra (2004), all in subduction zones. The longest earthquake ruptures on strike-slip faults, like the San Andreas Fault (1857, 1906), the North Anatolian Fault in Turkey (1939), and the Denali Fault in Alaska (2002), are about half to one third as long as the lengths along subducting plate margins, and those along normal faults are even shorter.
+
+The most important parameter controlling the maximum earthquake magnitude on a fault, however, is not the maximum available length, but the available width because the latter varies by a factor of 20. Along converging plate margins, the dip angle of the rupture plane is very shallow, typically about 10 degrees.[8] Thus, the width of the plane within the top brittle crust of the Earth can become 50–100 km (31–62 mi) (Japan, 2011; Alaska, 1964), making the most powerful earthquakes possible.
+
+Strike-slip faults tend to be oriented near vertically, resulting in an approximate width of 10 km (6.2 mi) within the brittle crust.[9] Thus, earthquakes with magnitudes much larger than 8 are not possible. Maximum magnitudes along many normal faults are even more limited because many of them are located along spreading centers, as in Iceland, where the thickness of the brittle layer is only about six kilometres (3.7 mi).[10][11]
+
+In addition, there exists a hierarchy of stress level in the three fault types. Thrust faults are generated by the highest, strike-slip by intermediate, and normal faults by the lowest stress levels.[12] This can easily be understood by considering the direction of the greatest principal stress, the direction of the force that "pushes" the rock mass during the faulting. In the case of normal faults, the rock mass is pushed down in a vertical direction, thus the pushing force (greatest principal stress) equals the weight of the rock mass itself. In the case of thrusting, the rock mass "escapes" in the direction of the least principal stress, namely upward, lifting the rock mass up, and thus, the overburden equals the least principal stress. Strike-slip faulting is intermediate between the other two types described above. This difference in stress regime in the three faulting environments can contribute to differences in stress drop during faulting, which contributes to differences in the radiated energy, regardless of fault dimensions.
+
+Where plate boundaries occur within the continental lithosphere, deformation is spread out over a much larger area than the plate boundary itself. In the case of the San Andreas fault continental transform, many earthquakes occur away from the plate boundary and are related to strains developed within the broader zone of deformation caused by major irregularities in the fault trace (e.g., the "Big bend" region). The Northridge earthquake was associated with movement on a blind thrust within such a zone. Another example is the strongly oblique convergent plate boundary between the Arabian and Eurasian plates where it runs through the northwestern part of the Zagros Mountains. The deformation associated with this plate boundary is partitioned into nearly pure thrust sense movements perpendicular to the boundary over a wide zone to the southwest and nearly pure strike-slip motion along the Main Recent Fault close to the actual plate boundary itself. This is demonstrated by earthquake focal mechanisms.[13]
+
+All tectonic plates have internal stress fields caused by their interactions with neighboring plates and sedimentary loading or unloading (e.g., deglaciation).[14] These stresses may be sufficient to cause failure along existing fault planes, giving rise to intraplate earthquakes.[15]
+
+The majority of tectonic earthquakes originate at the ring of fire in depths not exceeding tens of kilometers. Earthquakes occurring at a depth of less than 70 km (43 mi) are classified as "shallow-focus" earthquakes, while those with a focal-depth between 70 and 300 km (43 and 186 mi) are commonly termed "mid-focus" or "intermediate-depth" earthquakes. In subduction zones, where older and colder oceanic crust descends beneath another tectonic plate, deep-focus earthquakes may occur at much greater depths (ranging from 300 to 700 km (190 to 430 mi)).[16] These seismically active areas of subduction are known as Wadati–Benioff zones. Deep-focus earthquakes occur at a depth where the subducted lithosphere should no longer be brittle, due to the high temperature and pressure. A possible mechanism for the generation of deep-focus earthquakes is faulting caused by olivine undergoing a phase transition into a spinel structure.[17]
+
+Earthquakes often occur in volcanic regions and are caused there, both by tectonic faults and the movement of magma in volcanoes. Such earthquakes can serve as an early warning of volcanic eruptions, as during the 1980 eruption of Mount St. Helens.[18] Earthquake swarms can serve as markers for the location of the flowing magma throughout the volcanoes. These swarms can be recorded by seismometers and tiltmeters (a device that measures ground slope) and used as sensors to predict imminent or upcoming eruptions.[19]
+
+A tectonic earthquake begins by an initial rupture at a point on the fault surface, a process known as nucleation. The scale of the nucleation zone is uncertain, with some evidence, such as the rupture dimensions of the smallest earthquakes, suggesting that it is smaller than 100 m (330 ft) while other evidence, such as a slow component revealed by low-frequency spectra of some earthquakes, suggest that it is larger. The possibility that the nucleation involves some sort of preparation process is supported by the observation that about 40% of earthquakes are preceded by foreshocks. Once the rupture has initiated, it begins to propagate along the fault surface. The mechanics of this process are poorly understood, partly because it is difficult to recreate the high sliding velocities in a laboratory. Also the effects of strong ground motion make it very difficult to record information close to a nucleation zone.[20]
+
+Rupture propagation is generally modeled using a fracture mechanics approach, likening the rupture to a propagating mixed mode shear crack. The rupture velocity is a function of the fracture energy in the volume around the crack tip, increasing with decreasing fracture energy. The velocity of rupture propagation is orders of magnitude faster than the displacement velocity across the fault. Earthquake ruptures typically propagate at velocities that are in the range 70–90% of the S-wave velocity, which is independent of earthquake size. A small subset of earthquake ruptures appear to have propagated at speeds greater than the S-wave velocity. These supershear earthquakes have all been observed during large strike-slip events. The unusually wide zone of coseismic damage caused by the 2001 Kunlun earthquake has been attributed to the effects of the sonic boom developed in such earthquakes. Some earthquake ruptures travel at unusually low velocities and are referred to as slow earthquakes. A particularly dangerous form of slow earthquake is the tsunami earthquake, observed where the relatively low felt intensities, caused by the slow propagation speed of some great earthquakes, fail to alert the population of the neighboring coast, as in the 1896 Sanriku earthquake.[20]
+
+Tides may induce some seismicity. See tidal triggering of earthquakes for details.
+
+Most earthquakes form part of a sequence, related to each other in terms of location and time.[21] Most earthquake clusters consist of small tremors that cause little to no damage, but there is a theory that earthquakes can recur in a regular pattern.[22]
+
+An aftershock is an earthquake that occurs after a previous earthquake, the mainshock. An aftershock is in the same region of the main shock but always of a smaller magnitude. If an aftershock is larger than the main shock, the aftershock is redesignated as the main shock and the original main shock is redesignated as a foreshock. Aftershocks are formed as the crust around the displaced fault plane adjusts to the effects of the main shock.[21]
+
+Earthquake swarms are sequences of earthquakes striking in a specific area within a short period of time. They are different from earthquakes followed by a series of aftershocks by the fact that no single earthquake in the sequence is obviously the main shock, so none has a notable higher magnitude than another. An example of an earthquake swarm is the 2004 activity at Yellowstone National Park.[23] In August 2012, a swarm of earthquakes shook Southern California's Imperial Valley, showing the most recorded activity in the area since the 1970s.[24]
+
+Sometimes a series of earthquakes occur in what has been called an earthquake storm, where the earthquakes strike a fault in clusters, each triggered by the shaking or stress redistribution of the previous earthquakes. Similar to aftershocks but on adjacent segments of fault, these storms occur over the course of years, and with some of the later earthquakes as damaging as the early ones. Such a pattern was observed in the sequence of about a dozen earthquakes that struck the North Anatolian Fault in Turkey in the 20th century and has been inferred for older anomalous clusters of large earthquakes in the Middle East.[25][26]
+
+Quaking or shaking of the earth is a common phenomenon undoubtedly known to humans from earliest times. Prior to the development of strong-motion accelerometers that can measure peak ground speed and acceleration directly, the intensity of the earth-shaking was estimated on the basis of the observed effects, as categorized on various seismic intensity scales. Only in the last century has the source of such shaking been identified as ruptures in the Earth's crust, with the intensity of shaking at any locality dependent not only on the local ground conditions but also on the strength or magnitude of the rupture, and on its distance.[27]
+
+The first scale for measuring earthquake magnitudes was developed by Charles F. Richter in 1935. Subsequent scales (see seismic magnitude scales) have retained a key feature, where each unit represents a ten-fold difference in the amplitude of the ground shaking and a 32-fold difference in energy. Subsequent scales are also adjusted to have approximately the same numeric value within the limits of the scale.[28]
+
+Although the mass media commonly reports earthquake magnitudes as "Richter magnitude" or "Richter scale", standard practice by most seismological authorities is to express an earthquake's strength on the moment magnitude scale, which is based on the actual energy released by an earthquake.[29]
+
+It is estimated that around 500,000 earthquakes occur each year, detectable with current instrumentation. About 100,000 of these can be felt.[30][31] Minor earthquakes occur nearly constantly around the world in places like California and Alaska in the U.S., as well as in El Salvador, Mexico, Guatemala, Chile, Peru, Indonesia, Philippines, Iran, Pakistan, the Azores in Portugal, Turkey, New Zealand, Greece, Italy, India, Nepal and Japan.[32] Larger earthquakes occur less frequently, the relationship being exponential; for example, roughly ten times as many earthquakes larger than magnitude 4 occur in a particular time period than earthquakes larger than magnitude 5.[33] In the (low seismicity) United Kingdom, for example, it has been calculated that the average recurrences are:
+an earthquake of 3.7–4.6 every year, an earthquake of 4.7–5.5 every 10 years, and an earthquake of 5.6 or larger every 100 years.[34] This is an example of the Gutenberg–Richter law.
+
+The number of seismic stations has increased from about 350 in 1931 to many thousands today. As a result, many more earthquakes are reported than in the past, but this is because of the vast improvement in instrumentation, rather than an increase in the number of earthquakes. The United States Geological Survey estimates that, since 1900, there have been an average of 18 major earthquakes (magnitude 7.0–7.9) and one great earthquake (magnitude 8.0 or greater) per year, and that this average has been relatively stable.[36] In recent years, the number of major earthquakes per year has decreased, though this is probably a statistical fluctuation rather than a systematic trend.[37] More detailed statistics on the size and frequency of earthquakes is available from the United States Geological Survey (USGS).[38]
+A recent increase in the number of major earthquakes has been noted, which could be explained by a cyclical pattern of periods of intense tectonic activity, interspersed with longer periods of low intensity. However, accurate recordings of earthquakes only began in the early 1900s, so it is too early to categorically state that this is the case.[39]
+
+Most of the world's earthquakes (90%, and 81% of the largest) take place in the 40,000-kilometre-long (25,000 mi), horseshoe-shaped zone called the circum-Pacific seismic belt, known as the Pacific Ring of Fire, which for the most part bounds the Pacific Plate.[40][41] Massive earthquakes tend to occur along other plate boundaries too, such as along the Himalayan Mountains.[42]
+
+With the rapid growth of mega-cities such as Mexico City, Tokyo and Tehran in areas of high seismic risk, some seismologists are warning that a single quake may claim the lives of up to three million people.[43]
+
+While most earthquakes are caused by movement of the Earth's tectonic plates, human activity can also produce earthquakes. Activities both above ground and below may change the stresses and strains on the crust, including building reservoirs, extracting resources such as coal or oil, and injecting fluids underground for waste disposal or fracking.[44] Most of these earthquakes have small magnitudes. The 5.7 magnitude 2011 Oklahoma earthquake is thought to have been caused by disposing wastewater from oil production into injection wells,[45] and studies point to the state's oil industry as the cause of other earthquakes in the past century.[46] A Columbia University paper suggested that the 8.0 magnitude 2008 Sichuan earthquake was induced by loading from the Zipingpu Dam, though the link has not been conclusively proved.[47]
+
+The instrumental scales used to describe the size of an earthquake began with the Richter magnitude scale in the 1930s. It is a relatively simple measurement of an event's amplitude, and its use has become minimal in the 21st century. Seismic waves travel through the Earth's interior and can be recorded by seismometers at great distances. The surface wave magnitude was developed in the 1950s as a means to measure remote earthquakes and to improve the accuracy for larger events. The moment magnitude scale not only measures the amplitude of the shock but also takes into account the seismic moment (total rupture area, average slip of the fault, and rigidity of the rock). The Japan Meteorological Agency seismic intensity scale, the Medvedev–Sponheuer–Karnik scale, and the Mercalli intensity scale are based on the observed effects and are related to the intensity of shaking.
+
+Every tremor produces different types of seismic waves, which travel through rock with different velocities:
+
+Propagation velocity of the seismic waves through solid rock ranges from approx. 3 km/s (1.9 mi/s) up to 13 km/s (8.1 mi/s), depending on the density and elasticity of the medium. In the Earth's interior, the shock- or P-waves travel much faster than the S-waves (approx. relation 1.7:1). The differences in travel time from the epicenter to the observatory are a measure of the distance and can be used to image both sources of quakes and structures within the Earth. Also, the depth of the hypocenter can be computed roughly.
+
+In the upper crust, P-waves travel in the range 2–3 km (1.2–1.9 mi) per second (or lower) in soils and unconsolidated sediments, increasing to 3–6 km (1.9–3.7 mi) per second in solid rock. In the lower crust, they travel at about 6–7 km (3.7–4.3 mi) per second; the velocity increases within the deep mantle to about 13 km (8.1 mi) per second. The velocity of S-waves ranges from 2–3 km (1.2–1.9 mi) per second in light sediments and 4–5 km (2.5–3.1 mi) per second in the Earth's crust up to 7 km (4.3 mi) per second in the deep mantle. As a consequence, the first waves of a distant earthquake arrive at an observatory via the Earth's mantle.
+
+On average, the kilometer distance to the earthquake is the number of seconds between the P- and S-wave times 8.[48] Slight deviations are caused by inhomogeneities of subsurface structure. By such analyses of seismograms the Earth's core was located in 1913 by Beno Gutenberg.
+
+S-waves and later arriving surface waves do most of the damage compared to P-waves. P-waves squeeze and expand material in the same direction they are traveling, whereas S-waves shake the ground up and down and back and forth.[49]
+
+Earthquakes are not only categorized by their magnitude but also by the place where they occur. The world is divided into 754 Flinn–Engdahl regions (F-E regions), which are based on political and geographical boundaries as well as seismic activity. More active zones are divided into smaller F-E regions whereas less active zones belong to larger F-E regions.
+
+Standard reporting of earthquakes includes its magnitude, date and time of occurrence, geographic coordinates of its epicenter, depth of the epicenter, geographical region, distances to population centers, location uncertainty, a number of parameters that are included in USGS earthquake reports (number of stations reporting, number of observations, etc.), and a unique event ID.[50]
+
+Although relatively slow seismic waves have traditionally been used to detect earthquakes, scientists realized in 2016 that gravitational measurements could provide instantaneous detection of earthquakes, and confirmed this by analyzing gravitational records associated with the 2011 Tohoku-Oki ("Fukushima") earthquake.[51][52]
+
+The effects of earthquakes include, but are not limited to, the following:
+
+Shaking and ground rupture are the main effects created by earthquakes, principally resulting in more or less severe damage to buildings and other rigid structures. The severity of the local effects depends on the complex combination of the earthquake magnitude, the distance from the epicenter, and the local geological and geomorphological conditions, which may amplify or reduce wave propagation.[53] The ground-shaking is measured by ground acceleration.
+
+Specific local geological, geomorphological, and geostructural features can induce high levels of shaking on the ground surface even from low-intensity earthquakes. This effect is called site or local amplification. It is principally due to the transfer of the seismic motion from hard deep soils to soft superficial soils and to effects of seismic energy focalization owing to typical geometrical setting of the deposits.
+
+Ground rupture is a visible breaking and displacement of the Earth's surface along the trace of the fault, which may be of the order of several meters in the case of major earthquakes. Ground rupture is a major risk for large engineering structures such as dams, bridges, and nuclear power stations and requires careful mapping of existing faults to identify any that are likely to break the ground surface within the life of the structure.[54]
+
+Soil liquefaction occurs when, because of the shaking, water-saturated granular material (such as sand) temporarily loses its strength and transforms from a solid to a liquid. Soil liquefaction may cause rigid structures, like buildings and bridges, to tilt or sink into the liquefied deposits. For example, in the 1964 Alaska earthquake, soil liquefaction caused many buildings to sink into the ground, eventually collapsing upon themselves.[55]
+
+An earthquake may cause injury and loss of life, road and bridge damage, general property damage, and collapse or destabilization (potentially leading to future collapse) of buildings. The aftermath may bring disease, lack of basic necessities, mental consequences such as panic attacks, depression to survivors,[56] and higher insurance premiums.
+
+Earthquakes can produce slope instability leading to landslides, a major geological hazard. Landslide danger may persist while emergency personnel are attempting rescue.[57]
+
+Earthquakes can cause fires by damaging electrical power or gas lines. In the event of water mains rupturing and a loss of pressure, it may also become difficult to stop the spread of a fire once it has started. For example, more deaths in the 1906 San Francisco earthquake were caused by fire than by the earthquake itself.[58]
+
+Tsunamis are long-wavelength, long-period sea waves produced by the sudden or abrupt movement of large volumes of water—including when an earthquake occurs at sea. In the open ocean the distance between wave crests can surpass 100 kilometers (62 mi), and the wave periods can vary from five minutes to one hour. Such tsunamis travel 600–800 kilometers per hour (373–497 miles per hour), depending on water depth. Large waves produced by an earthquake or a submarine landslide can overrun nearby coastal areas in a matter of minutes. Tsunamis can also travel thousands of kilometers across open ocean and wreak destruction on far shores hours after the earthquake that generated them.[59]
+
+Ordinarily, subduction earthquakes under magnitude 7.5 do not cause tsunamis, although some instances of this have been recorded. Most destructive tsunamis are caused by earthquakes of magnitude 7.5 or more.[59]
+
+Floods may be secondary effects of earthquakes, if dams are damaged. Earthquakes may cause landslips to dam rivers, which collapse and cause floods.[60]
+
+The terrain below the Sarez Lake in Tajikistan is in danger of catastrophic flooding if the landslide dam formed by the earthquake, known as the Usoi Dam, were to fail during a future earthquake. Impact projections suggest the flood could affect roughly 5 million people.[61]
+
+One of the most devastating earthquakes in recorded history was the 1556 Shaanxi earthquake, which occurred on 23 January 1556 in Shaanxi province, China. More than 830,000 people died.[63] Most houses in the area were yaodongs—dwellings carved out of loess hillsides—and many victims were killed when these structures collapsed. The 1976 Tangshan earthquake, which killed between 240,000 and 655,000 people, was the deadliest of the 20th century.[64]
+
+The 1960 Chilean earthquake is the largest earthquake that has been measured on a seismograph, reaching 9.5 magnitude on 22 May 1960.[30][31] Its epicenter was near Cañete, Chile. The energy released was approximately twice that of the next most powerful earthquake, the Good Friday earthquake (March 27, 1964), which was centered in Prince William Sound, Alaska.[65][66] The ten largest recorded earthquakes have all been megathrust earthquakes; however, of these ten, only the 2004 Indian Ocean earthquake is simultaneously one of the deadliest earthquakes in history.
+
+Earthquakes that caused the greatest loss of life, while powerful, were deadly because of their proximity to either heavily populated areas or the ocean, where earthquakes often create tsunamis that can devastate communities thousands of kilometers away. Regions most at risk for great loss of life include those where earthquakes are relatively rare but powerful, and poor regions with lax, unenforced, or nonexistent seismic building codes.
+
+Earthquake prediction is a branch of the science of seismology concerned with the specification of the time, location, and magnitude of future earthquakes within stated limits.[67] Many methods have been developed for predicting the time and place in which earthquakes will occur. Despite considerable research efforts by seismologists, scientifically reproducible predictions cannot yet be made to a specific day or month.[68]
+
+While forecasting is usually considered to be a type of prediction, earthquake forecasting is often differentiated from earthquake prediction. Earthquake forecasting is concerned with the probabilistic assessment of general earthquake hazard, including the frequency and magnitude of damaging earthquakes in a given area over years or decades.[69] For well-understood faults the probability that a segment may rupture during the next few decades can be estimated.[70][71]
+
+Earthquake warning systems have been developed that can provide regional notification of an earthquake in progress, but before the ground surface has begun to move, potentially allowing people within the system's range to seek shelter before the earthquake's impact is felt.
+
+The objective of earthquake engineering is to foresee the impact of earthquakes on buildings and other structures and to design such structures to minimize the risk of damage. Existing structures can be modified by seismic retrofitting to improve their resistance to earthquakes. Earthquake insurance can provide building owners with financial protection against losses resulting from earthquakes Emergency management strategies can be employed by a government or organization to mitigate risks and prepare for consequences.
+
+Individuals can also take preparedness steps like securing water heaters and heavy items that could injure someone, locating shutoffs for utilities, and being educated about what to do when shaking starts. For areas near large bodies of water, earthquake preparedness encompasses the possibility of a tsunami caused by a large quake.
+
+From the lifetime of the Greek philosopher Anaxagoras in the 5th century BCE to the 14th century CE, earthquakes were usually attributed to "air (vapors) in the cavities of the Earth."[72] Thales of Miletus (625–547 BCE) was the only documented person who believed that earthquakes were caused by tension between the earth and water.[72] Other theories existed, including the Greek philosopher Anaxamines' (585–526 BCE) beliefs that short incline episodes of dryness and wetness caused seismic activity. The Greek philosopher Democritus (460–371 BCE) blamed water in general for earthquakes.[72] Pliny the Elder called earthquakes "underground thunderstorms".[72]
+
+In recent studies, geologists claim that global warming is one of the reasons for increased seismic activity. According to these studies, melting glaciers and rising sea levels disturb the balance of pressure on Earth's tectonic plates, thus causing an increase in the frequency and intensity of earthquakes.[73]
+
+In Norse mythology, earthquakes were explained as the violent struggling of the god Loki. When Loki, god of mischief and strife, murdered Baldr, god of beauty and light, he was punished by being bound in a cave with a poisonous serpent placed above his head dripping venom. Loki's wife Sigyn stood by him with a bowl to catch the poison, but whenever she had to empty the bowl the poison dripped on Loki's face, forcing him to jerk his head away and thrash against his bonds, which caused the earth to tremble.[74]
+
+In Greek mythology, Poseidon was the cause and god of earthquakes. When he was in a bad mood, he struck the ground with a trident, causing earthquakes and other calamities. He also used earthquakes to punish and inflict fear upon people as revenge.[75]
+
+In Japanese mythology, Namazu (鯰) is a giant catfish who causes earthquakes. Namazu lives in the mud beneath the earth, and is guarded by the god Kashima who restrains the fish with a stone. When Kashima lets his guard fall, Namazu thrashes about, causing violent earthquakes.[76]
+
+In modern popular culture, the portrayal of earthquakes is shaped by the memory of great cities laid waste, such as Kobe in 1995 or San Francisco in 1906.[77] Fictional earthquakes tend to strike suddenly and without warning.[77] For this reason, stories about earthquakes generally begin with the disaster and focus on its immediate aftermath, as in Short Walk to Daylight (1972), The Ragged Edge (1968) or Aftershock: Earthquake in New York (1999).[77] A notable example is Heinrich von Kleist's classic novella, The Earthquake in Chile, which describes the destruction of Santiago in 1647. Haruki Murakami's short fiction collection After the Quake depicts the consequences of the Kobe earthquake of 1995.
+
+The most popular single earthquake in fiction is the hypothetical "Big One" expected of California's San Andreas Fault someday, as depicted in the novels Richter 10 (1996), Goodbye California (1977), 2012 (2009) and San Andreas (2015) among other works.[77] Jacob M. Appel's widely anthologized short story, A Comparative Seismology, features a con artist who convinces an elderly woman that an apocalyptic earthquake is imminent.[78]
+
+Contemporary depictions of earthquakes in film are variable in the manner in which they reflect human psychological reactions to the actual trauma that can be caused to directly afflicted families and their loved ones.[79] Disaster mental health response research emphasizes the need to be aware of the different roles of loss of family and key community members, loss of home and familiar surroundings, loss of essential supplies and services to maintain survival.[80][81] Particularly for children, the clear availability of caregiving adults who are able to protect, nourish, and clothe them in the aftermath of the earthquake, and to help them make sense of what has befallen them has been shown even more important to their emotional and physical health than the simple giving of provisions.[82] As was observed after other disasters involving destruction and loss of life and their media depictions, recently observed in the 2010 Haiti earthquake, it is also important not to pathologize the reactions to loss and displacement or disruption of governmental administration and services, but rather to validate these reactions, to support constructive problem-solving and reflection as to how one might improve the conditions of those affected.[83]

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+Epiphany may refer to:

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+
+
+A bishop is an ordained, consecrated, or appointed member of the Christian clergy who is generally entrusted with a position of authority and oversight.
+
+Within the Catholic, Eastern Orthodox, Oriental Orthodox, Moravian, Anglican, Old Catholic and Independent Catholic churches, as well as the Assyrian Church of the East, bishops claim apostolic succession, a direct historical lineage dating back to the original Twelve Apostles. Within these churches, bishops are seen as those who possess the full priesthood and can ordain clergy, including other bishops. Some Protestant churches, including the Lutheran, Anglican and Methodist churches, have bishops serving similar functions as well, though not always understood to be within apostolic succession in the same way. A person ordained as a deacon, priest, and then bishop is understood to hold the fullness of the (ministerial) priesthood, given responsibility by Christ to govern, teach, and sanctify the Body of Christ. Priests, deacons and lay ministers co-operate and assist their bishops in pastoral ministry.
+
+The English term bishop derives from the Greek word ἐπίσκοπος epískopos, meaning "overseer" in Greek, the early language of the Christian Church. In the early Christian era the term was not always clearly distinguished from  presbýteros (literally: "elder" or "senior", origin of the modern English word "priest"), but is used in the sense of the order or office of bishop, distinct from that of presbyter, in the writings attributed to Ignatius of Antioch.[1] (died c. 110).
+
+The earliest organization of the Church in Jerusalem was, according to most scholars, similar to that of Jewish synagogues, but it had a council or college of ordained presbyters (Ancient Greek: πρεσβύτεροι elders). In Acts 11:30 and Acts 15:22, we see a collegiate system of government in Jerusalem chaired by James the Just, according to tradition the first bishop of the city. In Acts 14:23, the Apostle Paul ordains presbyters in churches in Anatolia.[2]
+
+Often, the word presbyter was not yet distinguished from overseer (Ancient Greek: ἐπίσκοπος episkopos, later used exclusively to mean bishop), as in Acts 20:17, Titus 1:5–7 and 1 Peter 5:1.[a][b] The earliest writings of the Apostolic Fathers, the Didache and the First Epistle of Clement, for example, show the church used two terms for local church offices—presbyters (seen by many as an interchangeable term with episcopos or overseer) and deacon.
+
+In Timothy and Titus in the New Testament a more clearly defined episcopate can be seen. We are told that Paul had left Timothy in Ephesus and Titus in Crete to oversee the local church.[6][7] Paul commands Titus to ordain presbyters/bishops and to exercise general oversight.
+
+Early sources are unclear but various groups of Christian communities may have had the bishop surrounded by a group or college functioning as leaders of the local churches.[8][9] Eventually the head or "monarchic" bishop came to rule more clearly,[10] and all local churches would eventually follow the example of the other churches and structure themselves after the model of the others with the one bishop in clearer charge,[8] though the role of the body of presbyters remained important.[10]
+
+Eventually, as Christendom grew, bishops no longer directly served individual congregations. Instead, the Metropolitan bishop (the bishop in a large city) appointed priests to minister each congregation, acting as the bishop's delegate.
+
+Around the end of the 1st century, the church's organization became clearer in historical documents. In the works of the Apostolic Fathers, and Ignatius of Antioch in particular, the role of the episkopos, or bishop, became more important or, rather, already was very important and being clearly defined. While Ignatius of Antioch offers the earliest clear description of monarchial bishops (a single bishop over all house churches in a city) he is an advocate of monepiscopal structure rather than describing an accepted reality. To the bishops and house churches to which he writes, he offers strategies on how to pressure house churches who don't recognize the bishop into compliance. Other contemporary Christian writers do not describe monarchial bishops, either continuing to equate them with the presbyters or speaking of episkopoi (bishops, plural) in a city.
+
+"Blessed be God, who has granted unto you, who are yourselves so excellent, to obtain such an excellent bishop." — Epistle of Ignatius to the Ephesians 1:1  [11] 
+"and that, being subject to the bishop and the presbytery, ye may in all respects be sanctified." — Epistle of Ignatius to the Ephesians 2:1
+[12] 
+"For your justly renowned presbytery, worthy of God, is fitted as exactly to the bishop as the strings are to the harp." — Epistle of Ignatius to the Ephesians 4:1 [13] 
+"Do ye, beloved, be careful to be subject to the bishop, and the presbyters and the deacons." — Epistle of Ignatius to the Ephesians 5:1 [13] 
+"Plainly therefore we ought to regard the bishop as the Lord Himself" — Epistle of Ignatius to the Ephesians 6:1.
+"your godly bishop" — Epistle of Ignatius to the Magnesians 2:1.
+"the bishop presiding after the likeness of God and the presbyters after the likeness of the council of the Apostles, with the deacons also who are most dear to me, having been entrusted with the diaconate of Jesus Christ" — Epistle of Ignatius to the Magnesians 6:1.
+"Therefore as the Lord did nothing without the Father, [being united with Him], either by Himself or by the Apostles, so neither do ye anything without the bishop and the presbyters." — Epistle of Ignatius to the Magnesians 7:1.
+"Be obedient to the bishop and to one another, as Jesus Christ was to the Father [according to the flesh], and as the Apostles were to Christ and to the Father, that there may be union both of flesh and of spirit." — Epistle of Ignatius to the Magnesians 13:2.
+"In like manner let all men respect the deacons as Jesus Christ, even as they should respect the bishop as being a type of the Father and the presbyters as the council of God and as the college of Apostles. Apart from these there is not even the name of a church." — Epistle of Ignatius to the Trallesians 3:1.
+"follow your bishop, as Jesus Christ followed the Father, and the presbytery as the Apostles; and to the deacons pay respect, as to God's commandment" — Epistle of Ignatius to the Smyrnans 8:1.
+"He that honoureth the bishop is honoured of God; he that doeth aught without the knowledge of the bishop rendereth service to the devil" — Epistle of Ignatius to the Smyrnans 9:1.
+
+As the Church continued to expand, new churches in important cities gained their own bishop. Churches in the regions outside an important city were served by Chorbishop, an official rank of bishops. However, soon, presbyters and deacons were sent from bishop of a city church. Gradually priests replaced the chorbishops. Thus, in time, the bishop changed from being the leader of a single church confined to an urban area to being the leader of the churches of a given geographical area.
+
+Clement of Alexandria (end of the 2nd century) writes about the ordination of a certain Zachæus as bishop by the imposition of Simon Peter Bar-Jonah's hands. The words bishop and ordination are used in their technical meaning by the same Clement of Alexandria.[14] The bishops in the 2nd century are defined also as the only clergy to whom the ordination to priesthood (presbyterate) and diaconate is entrusted: "a priest (presbyter) lays on hands, but does not ordain." (cheirothetei ou cheirotonei[15])
+
+At the beginning of the 3rd century, Hippolytus of Rome describes another feature of the ministry of a bishop, which is that of the "Spiritum primatus sacerdotii habere potestatem dimittere peccata": the primate of sacrificial priesthood and the power to forgive sins.[16]
+
+The efficient organization of the Roman Empire became the template for the organisation of the church in the 4th century, particularly after Constantine's Edict of Milan. As the church moved from the shadows of privacy into the public forum it acquired land for churches, burials and clergy. In 391, Theodosius I decreed that any land that had been confiscated from the church by Roman authorities be returned.
+
+The most usual term for the geographic area of a bishop's authority and ministry, the diocese, began as part of the structure of the Roman Empire under Diocletian. As Roman authority began to fail in the western portion of the empire, the church took over much of the civil administration. This can be clearly seen in the ministry of two popes: Pope Leo I in the 5th century, and Pope Gregory I in the 6th century. Both of these men were statesmen and public administrators in addition to their role as Christian pastors, teachers and leaders. In the Eastern churches, latifundia entailed to a bishop's see were much less common, the state power did not collapse the way it did in the West, and thus the tendency of bishops acquiring civil power was much weaker than in the West. However, the role of Western bishops as civil authorities, often called prince bishops, continued throughout much of the Middle Ages.
+
+As well as being archchancellors of the Holy Roman Empire after the 9th century, bishops generally served as chancellors to medieval monarchs, acting as head of the justiciary and chief chaplain. The Lord Chancellor of England was almost always a bishop up until the dismissal of Cardinal Thomas Wolsey by Henry VIII. Similarly, the position of Kanclerz in the Polish kingdom was always held by a bishop until the 16th century. And today, the principality of Andorra is headed by two co-princes, one of whom is a Catholic bishop (and the other, the President of France).
+
+In France before the French Revolution, representatives of the clergy — in practice, bishops and abbots of the largest monasteries — comprised the First Estate of the Estates-General, until their role was abolished during the French Revolution.
+
+In the 21st century, the more senior bishops of the Church of England continue to sit in the House of Lords of the Parliament of the United Kingdom, as representatives of the established church, and are known as Lords Spiritual. The Bishop of Sodor and Man, whose diocese lies outside the United Kingdom, is an ex officio member of the Legislative Council of the Isle of Man. In the past, the Bishop of Durham, known as a prince bishop, had extensive viceregal powers within his northern diocese — the power to mint money, collect taxes and raise an army to defend against the Scots.
+
+Eastern Orthodox bishops, along with all other members of the clergy, are canonically forbidden to hold political office. Occasional exceptions to this rule are tolerated when the alternative is political chaos. In the Ottoman Empire, the Patriarch of Constantinople, for example, had de facto administrative, fiscal, cultural and legal jurisdiction, as well as spiritual, over all the Christians of the empire. More recently, Archbishop Makarios III of Cyprus, served as President of the Republic of Cyprus from 1960 to 1977.
+
+In 2001, Peter Hollingworth, AC, OBE – then the Anglican Archbishop of Brisbane – was controversially appointed Governor-General of Australia. Although Hollingworth gave up his episcopal position to accept the appointment, it still attracted considerable opposition in a country which maintains a formal separation between Church and State.
+
+During the period of the English Civil War, the role of bishops as wielders of political power and as upholders of the established church became a matter of heated political controversy. Indeed, Presbyterianism was the polity of most Reformed Churches in Europe, and had been favored by many in England since the English Reformation.  Since in the primitive church the offices of presbyter and episkopos were identical, many Puritans held that this was the only form of government the church should have. The Anglican divine, Richard Hooker, objected to this claim in his famous work Of the Laws of Ecclesiastic Polity while, at the same time, defending Presbyterian ordination as valid (in particular Calvin's ordination of Beza). This was the official stance of the English Church until the Commonwealth, during which time, the views of Presbyterians and Independents (Congregationalists) were more freely expressed and practiced.
+
+Bishops form the leadership in the Catholic Church, the Eastern Orthodox Church, the Oriental Orthodox Churches, the Anglican Communion, the Lutheran Church, the Independent Catholic Churches, the Independent Anglican Churches, and certain other, smaller, denominations.
+
+The traditional role of a bishop is as pastor of a diocese (also called a bishopric, synod, eparchy or see), and so to serve as a "diocesan bishop," or "eparch" as it is called in many Eastern Christian churches. Dioceses vary considerably in size, geographically and population-wise. Some dioceses around the Mediterranean Sea which were Christianised early are rather compact, whereas dioceses in areas of rapid modern growth in Christian commitment—as in some parts of Sub-Saharan Africa, South America and the Far East—are much larger and more populous.
+
+As well as traditional diocesan bishops, many churches have a well-developed structure of church leadership that involves a number of layers of authority and responsibility.
+
+In Catholicism, Eastern Orthodoxy, Oriental Orthodoxy, and Anglicanism, only a bishop can ordain other bishops, priests, and deacons.
+
+In the Eastern liturgical tradition, a priest can celebrate the Divine Liturgy only with the blessing of a bishop. In Byzantine usage, an antimension signed by the bishop is kept on the altar partly as a reminder of whose altar it is and under whose omophorion the priest at a local parish is serving. In Syriac Church usage, a consecrated wooden block called a thabilitho is kept for the same reasons.
+
+The pope, in addition to being the Bishop of Rome and spiritual head of the Catholic Church, is also the Patriarch of the Latin Rite. Each bishop within the Latin Rite is answerable directly to the Pope and not any other bishop except to metropolitans in certain oversight instances. The pope previously used the title Patriarch of the West, but this title was dropped from use in 2006[18] a move which caused some concern within the Eastern Orthodox Communion as, to them, it implied wider papal jurisdiction.[19]
+
+In Catholic, Eastern Orthodox, Oriental Orthodox and Anglican cathedrals there is a special chair set aside for the exclusive use of the bishop. This is the bishop's cathedra and is often called the throne. In some Christian denominations, for example, the Anglican Communion, parish churches may maintain a chair for the use of the bishop when he visits; this is to signify the parish's union with the bishop.
+
+The bishop is the ordinary minister of the sacrament of confirmation in the Latin Rite Catholic Church, and in the Anglican and Old Catholic communion only a bishop may administer this sacrament. However, in the Byzantine and other Eastern rites, whether Eastern or Oriental Orthodox or Eastern Catholic, chrismation is done immediately after baptism, and thus the priest is the one who confirms, using chrism blessed by a bishop.[20]
+
+Bishops in all of these communions are ordained by other bishops through the laying on of hands. While traditional teaching maintains that any bishop with apostolic succession can validly perform the ordination of another bishop, some churches require two or three bishops participate, either to ensure sacramental validity or to conform with church law. Catholic doctrine holds that one bishop can validly ordain another (priest) as a bishop. Though a minimum of three bishops participating is desirable (there are usually several more) in order to demonstrate collegiality, canonically only one bishop is necessary. The practice of only one bishop ordaining was normal in countries where the Church was persecuted under Communist rule.
+The title of archbishop or metropolitan may be granted to a senior bishop, usually one who is in charge of a large ecclesiastical jurisdiction. He may, or may not, have provincial oversight of suffragan bishops and may possibly have auxiliary bishops assisting him.
+Ordination of a bishop, and thus continuation of apostolic succession, takes place through a ritual centred on the imposition of hands and prayer.
+Apart from the ordination, which is always done by other bishops, there are different methods as to the actual selection of a candidate for ordination as bishop. In the Catholic Church the Congregation for Bishops generally oversees the selection of new bishops with the approval of the pope. The papal nuncio usually solicits names from the bishops of a country, consults with priests and leading members of a laity, and then selects three to be forwarded to the Holy See. In Europe, some cathedral chapters have duties to elect bishops.  The Eastern Catholic churches generally elect their own bishops. Most Eastern Orthodox churches allow varying amounts of formalised laity or lower clergy influence on the choice of bishops. This also applies in those Eastern churches which are in union with the pope, though it is required that he give assent.
+
+Catholic, Eastern Orthodox, Oriental Orthodox, Anglican, Old Catholic and some Lutheran bishops claim to be part of the continuous sequence of ordained bishops since the days of the apostles referred to as apostolic succession. Since Pope Leo XIII issued the bull Apostolicae curae in 1896, the Catholic Church has insisted that Anglican orders are invalid because of changes in the Anglican ordination rites of the 16th century and divergence in understanding of the theology of priesthood, episcopacy and Eucharist. However, since the 1930s, Utrecht Old Catholic bishops (recognised by the Holy See as validily ordained) have sometimes taken part in the ordination of Anglican bishops. According to the writer Timothy Dufort, by 1969, all Church of England bishops had acquired Old Catholic lines of apostolic succession recognised by the Holy See.[21] This development has muddied the waters somewhat as it could be argued that the strain of apostolic succession has been re-introduced into Anglicanism, at least within the Church of England.
+
+The Catholic Church does recognise as valid (though illicit) ordinations done by breakaway Catholic, Old Catholic or Oriental bishops, and groups descended from them; it also regards as both valid and licit those ordinations done by bishops of the Eastern churches,[c] so long as those receiving the ordination conform to other canonical requirements (for example, is an adult male) and an eastern orthodox rite of episcopal ordination, expressing the proper functions and sacramental status of a bishop, is used; this has given rise to the phenomenon of episcopi vagantes (for example, clergy of the Independent Catholic groups which claim apostolic succession, though this claim is rejected by both Catholicism and Eastern Orthodoxy).
+
+The Eastern Orthodox Churches would not accept the validity of any ordinations performed by the Independent Catholic groups, as Eastern Orthodoxy considers to be spurious any consecration outside the Church as a whole. Eastern Orthodoxy considers apostolic succession to exist only within the Universal Church, and not through any authority held by individual bishops; thus, if a bishop ordains someone to serve outside the (Eastern Orthodox) Church, the ceremony is ineffectual, and no ordination has taken place regardless of the ritual used or the ordaining prelate's position within the Eastern Orthodox Churches.
+
+The position of the Catholic Church is slightly different. Whilst it does recognise the validity of the orders of certain groups which separated from communion with Holy See. The Holy See accepts as valid the ordinations of the Old Catholics in communion with Utrecht, as well as the Polish National Catholic Church (which received its orders directly from Utrecht, and was—until recently—part of that communion); but Catholicism does not recognise the orders of any group whose teaching is at variance with what they consider the core tenets of Christianity; this is the case even though the clergy of the Independent Catholic groups may use the proper ordination ritual. There are also other reasons why the Holy See does not recognise the validity of the orders of the Independent clergy:
+
+Whilst members of the Independent Catholic movement take seriously the issue of valid orders, it is highly significant that the relevant Vatican Congregations tend not to respond to petitions from Independent Catholic bishops and clergy who seek to be received into communion with the Holy See, hoping to continue in some sacramental role. In those instances where the pope does grant reconciliation, those deemed to be clerics within the Independent Old Catholic movement are invariably admitted as laity and not priests or bishops.
+
+There is a mutual recognition of the validity of orders amongst Catholic, Eastern Orthodox, Old Catholic, Oriental Orthodox and Assyrian Church of the East churches.[22]
+
+Some provinces of the Anglican Communion have begun ordaining women as bishops in recent decades – for example, England, Ireland, Scotland, Wales, the United States, Australia, New Zealand, Canada and Cuba. The first woman to be consecrated a bishop within Anglicanism was Barbara Harris, who was ordained in the United States in 1989. In 2006, Katharine Jefferts Schori, the Episcopal Bishop of Nevada, became the first woman to become the presiding bishop of the Episcopal Church.
+
+In the Evangelical Lutheran Church in America (ELCA) and the Evangelical Lutheran Church in Canada (ELCIC), the largest Lutheran Church bodies in the United States and Canada, respectively, and roughly based on the Nordic Lutheran state churches (similar to that of the Church of England), bishops are elected by Synod Assemblies, consisting of both lay members and clergy, for a term of six years, which can be renewed, depending upon the local synod's "constitution" (which is mirrored on either the ELCA or ELCIC's national constitution). Since the implementation of concordats between the ELCA and the Episcopal Church of the United States and the ELCIC and the Anglican Church of Canada, all bishops, including the presiding bishop (ELCA) or the national bishop (ELCIC), have been consecrated using the historic succession, with at least one Anglican bishop serving as co-consecrator.[23][24]
+
+Since going into ecumenical communion with their respective Anglican body, bishops in the ELCA or the ELCIC not only approve the "rostering" of all ordained pastors, diaconal ministers, and associates in ministry, but they serve as the principal celebrant of all pastoral ordination and installation ceremonies, diaconal consecration ceremonies, as well as serving as the "chief pastor" of the local synod, upholding the teachings of Martin Luther as well as the documentations of the Ninety-Five Theses and the Augsburg Confession. Unlike their counterparts in the United Methodist Church, ELCA and ELCIC synod bishops do not appoint pastors to local congregations (pastors, like their counterparts in the Episcopal Church, are called by local congregations). The presiding bishop of the ELCA and the national bishop of the ELCIC, the national bishops of their respective bodies, are elected for a single 6-year term and may be elected to an additional term.
+
+Although ELCA agreed with the Episcopal Church to limit ordination to the bishop "ordinarily", ELCA pastor-ordinators are given permission to perform the rites in "extraordinary" circumstance. In practice, "extraordinary" circumstance have included disagreeing with Episcopalian views of the episcopate, and as a result, ELCA pastors ordained by other pastors are not permitted to be deployed to Episcopal Churches (they can, however, serve in Presbyterian Church USA, United Methodist Church, Reformed Church in America, and Moravian Church congregations, as the ELCA is in full communion with these denominations). The Lutheran Church–Missouri Synod (LCMS) and the Wisconsin Evangelical Lutheran Synod (WELS), the second and third largest Lutheran bodies in the United States and the two largest Confessional Lutheran bodies in North America, do not follow an episcopal form of governance, settling instead on a form of quasi-congregationalism patterned off what they believe to be the practice of the early church. The second largest of the three predecessor bodies of the ELCA, the American Lutheran Church, was a congregationalist body, with national and synod presidents before they were re-titled as bishops (borrowing from the Lutheran churches in Germany) in the 1980s. With regard to ecclesial discipline and oversight, national and synod presidents typically function similarly to bishops in episcopal bodies.[25]
+
+In the African Methodist Episcopal Church, "Bishops are the Chief Officers of the Connectional Organization. They are elected for life by a majority vote of the General Conference which meets every four years."[26]
+
+In the Christian Methodist Episcopal Church in the United States, bishops are administrative superintendents of the church; they are elected by "delegate" votes for as many years deemed until the age of 74, then he/she must retire. Among their duties, are responsibility for appointing clergy to serve local churches as pastor, for performing ordinations, and for safeguarding the doctrine and discipline of the Church. The General Conference, a meeting every four years, has an equal number of clergy and lay delegates. In each Annual Conference, CME bishops serve for four-year terms. CME Church bishops may be male or female.
+
+In the United Methodist Church (the largest branch of Methodism in the world) bishops serve as administrative and pastoral superintendents of the church. They are elected for life from among the ordained elders (presbyters) by vote of the delegates in regional (called jurisdictional) conferences, and are consecrated by the other bishops present at the conference through the laying on of hands. In the United Methodist Church bishops remain members of the "Order of Elders" while being consecrated to the "Office of the Episcopacy". Within the United Methodist Church only bishops are empowered to consecrate bishops and ordain clergy. Among their most critical duties is the ordination and appointment of clergy to serve local churches as pastor, presiding at sessions of the Annual, Jurisdictional, and General Conferences, providing pastoral ministry for the clergy under their charge, and safeguarding the doctrine and discipline of the Church. Furthermore, individual bishops, or the Council of Bishops as a whole, often serve a prophetic role, making statements on important social issues and setting forth a vision for the denomination, though they have no legislative authority of their own. In all of these areas, bishops of the United Methodist Church function very much in the historic meaning of the term. According to the Book of Discipline of the United Methodist Church, a bishop's responsibilities are
+
+Leadership.—Spiritual and Temporal—
+
+Presidential Duties.—1. To preside in the General, Jurisdictional, Central, and Annual Conferences. 2. To form the districts after consultation with the district superintendents and after the number of the same has been determined by vote of the Annual Conference. 3. To appoint the district superintendents annually (¶¶ 517–518). 4. To consecrate bishops, to ordain elders and deacons, to consecrate diaconal ministers, to commission deaconesses and home missionaries, and to see that the names of the persons commissioned and consecrated are entered on the journals of the conference and that proper credentials are funised to these persons.
+
+Working with Ministers.—1. To make and fix the appointments in the Annual Conferences, Provisional Annual Conferences, and Missions as the Discipline may direct (¶¶ 529–533).
+
+2. To divide or to unite a circuit(s), stations(s), or mission(s) as judged necessary for missionary strategy and then to make appropriate appointments. 3. To read the appointments of deaconesses, diaconal ministers, lay persons in service under the World Division of the General Board of Global Ministries, and home missionaries. 4. To fix the Charge Conference membership of all ordained ministers appointed to ministries other than the local church in keeping with ¶443.3. 5. To transfer, upon the request of the receiving bishop, ministerial member(s) of one Annual Conference to another, provided said member(s) agrees to transfer; and to send immediately to the secretaries of both conferences involved, to the conference Boards of Ordained Ministry, and to the clearing house of the General Board of Pensions written notices of the transfer of members and of their standing in the course of study if they are undergraduates.[27]
+
+In each Annual Conference, United Methodist bishops serve for four-year terms, and may serve up to three terms before either retirement or appointment to a new Conference. United Methodist bishops may be male or female, with Marjorie Matthews being the first woman to be consecrated a bishop in 1980.
+
+The collegial expression of episcopal leadership in the United Methodist Church is known as the Council of Bishops. The Council of Bishops speaks to the Church and through the Church into the world and gives leadership in the quest for Christian unity and interreligious relationships.[27] The Conference of Methodist Bishops includes the United Methodist Council of Bishops plus bishops from affiliated autonomous Methodist or United Churches.
+
+John Wesley consecrated Thomas Coke a "General Superintendent," and directed that Francis Asbury also be consecrated for the United States of America in 1784, where the Methodist Episcopal Church first became a separate denomination apart from the Church of England. Coke soon returned to England, but Asbury was the primary builder of the new church. At first he did not call himself bishop, but eventually submitted to the usage by the denomination.
+
+Notable bishops in United Methodist history include Coke, Asbury, Richard Whatcoat, Philip William Otterbein, Martin Boehm, Jacob Albright, John Seybert, Matthew Simpson, John S. Stamm, William Ragsdale Cannon, Marjorie Matthews, Leontine T. Kelly, William B. Oden, Ntambo Nkulu Ntanda, Joseph Sprague, William Henry Willimon, and Thomas Bickerton.
+
+In The Church of Jesus Christ of Latter-day Saints, the Bishop is the leader of a local congregation, called a ward. As with most LDS priesthood holders, the bishop is a part-time lay minister and earns a living through other employment; in all cases, he is a married man. As such, it is his duty to preside at services, call local leaders, and judge the worthiness of members for service. The bishop does not deliver sermons at every service (generally asking members to do so), but is expected to be a spiritual guide for his congregation. It is therefore believed that he has both the right and ability to receive divine inspiration (through the Holy Spirit) for the ward under his direction. Because it is a part-time position, all able members are expected to assist in the management of the ward by holding delegated lay positions (for example, women's and youth leaders, teachers) referred to as callings. Although members are asked to confess serious sins to him, unlike the Catholic Church, he is not the instrument of divine forgiveness, but merely a guide through the repentance process (and a judge in case transgressions warrant excommunication or other official discipline). The bishop is also responsible for the physical welfare of the ward, and thus collects tithing and fast offerings and distributes financial assistance where needed.
+
+A bishop is the president of the Aaronic priesthood in his ward (and is thus a form of Mormon Kohen; in fact, a literal descendant of Aaron has "legal right" to act as a bishop[28] after being found worthy and ordained by the First Presidency[29]). In the absence of a literal descendant of Aaron, a high priest in the Melchizedek priesthood is called to be a bishop.[29] Each bishop is selected from resident members of the ward by the stake presidency with approval of the First Presidency, and chooses two counselors to form a bishopric. In special circumstances (such as a ward consisting entirely of young university students), a bishop may be chosen from outside the ward. A bishop is typically released after about five years and a new bishop is called to the position. Although the former bishop is released from his duties, he continues to hold the Aaronic priesthood office of bishop. Church members frequently refer to a former bishop as "Bishop" as a sign of respect and affection.
+
+Latter-day Saint bishops do not wear any special clothing or insignia the way clergy in many other churches do, but are expected to dress and groom themselves neatly and conservatively per their local culture, especially when performing official duties. Bishops (as well as other members of the priesthood) can trace their line of authority back to Joseph Smith, who, according to church doctrine, was ordained to lead the Church in modern times by the ancient apostles Peter, James, and John, who were ordained to lead the Church by Jesus Christ.[30]
+
+At the global level, the presiding bishop oversees the temporal affairs (buildings, properties, commercial corporations, and so on) of the worldwide Church, including the Church's massive global humanitarian aid and social welfare programs. The presiding bishop has two counselors; the three together form the presiding bishopric.[31] As opposed to ward bishoprics, where the counselors do not hold the office of bishop, all three men in the presiding bishopric hold the office of bishop, and thus the counselors, as with the presiding bishop, are formally referred to as "Bishop".[32]
+
+The New Apostolic Church (NAC) knows three classes of ministries: Deacons, Priests and Apostles. The Apostles, who are all included in the apostolate with the Chief Apostle as head, are the highest ministries.
+
+Of the several kinds of priest....ministries, the bishop is the highest. Nearly all bishops are set in line directly from the chief apostle. They support and help their superior apostle.
+
+In the Church of God in Christ (COGIC), the ecclesiastical structure is composed of large dioceses that are called "jurisdictions" within COGIC, each under the authority of a bishop, sometimes called "state bishops". They can either be made up of large geographical regions of churches or churches that are grouped and organized together as their own separate jurisdictions because of similar affiliations, regardless of geographical location or dispersion. Each state in the U.S. has at least one jurisdiction while others may have several more, and each jurisdiction is usually composed of between 30 and 100 churches. Each jurisdiction is then broken down into several districts, which are smaller groups of churches (either grouped by geographical situation or by similar affiliations) which are each under the authority of District Superintendents who answer to the authority of their jurisdictional/state bishop. There are currently over 170 jurisdictions in the United States, and over 30 jurisdictions in other countries. The bishops of each jurisdiction, according to the COGIC Manual, are considered to be the modern day equivalent in the church of the early apostles and overseers of the New Testament church, and as the highest ranking clergymen in the COGIC, they are tasked with the responsibilities of being the head overseers of all religious, civil, and economic ministries and protocol for the church denomination.[33] They also have the authority to appoint and ordain local pastors, elders, ministers, and reverends within the denomination. The bishops of the COGIC denomination are all collectively called "The Board of Bishops."[34] From the Board of Bishops, and the General Assembly of the COGIC, the body of the church composed of clergy and lay delegates that are responsible for making and enforcing the bylaws of the denomination, every four years, twelve bishops from the COGIC are elected as "The General Board" of the church, who work alongside the delegates of the General Assembly and Board of Bishops to provide administration over the denomination as the church's head executive leaders.[35] One of twelve bishops of the General Board is also elected the "presiding bishop" of the church, and two others are appointed by the presiding bishop himself, as his first and second assistant presiding bishops.
+
+Bishops in the Church of God in Christ usually wear black clergy suits which consist of a black suit blazer, black pants, a purple or scarlet clergy shirt and a white clerical collar, which is usually referred to as "Class B Civic attire." Bishops in COGIC also typically wear the Anglican Choir Dress style vestments of a long purple or scarlet chimere, cuffs, and tippet worn over a long white rochet, and a gold pectoral cross worn around the neck with the tippet. This is usually referred to as "Class A Ceremonial attire". The bishops of COGIC alternate between Class A Ceremonial attire and Class B Civic attire depending on the protocol of the religious services and other events they have to attend.[34][33]
+
+In the polity of the Church of God (Cleveland, Tennessee), the international leader is the presiding bishop, and the members of the executive committee are executive bishops. Collectively, they supervise and appoint national and state leaders across the world. Leaders of individual states and regions are administrative bishops, who have jurisdiction over local churches in their respective states and are vested with appointment authority for local pastorates. All ministers are credentialed at one of three levels of licensure, the most senior of which is the rank of ordained bishop. To be eligible to serve in state, national, or international positions of authority, a minister must hold the rank of ordained bishop.
+
+In 2002, the general convention of the Pentecostal Church of God came to a consensus to change the title of their overseer from general superintendent to bishop. The change was brought on because internationally, the term bishop is more commonly related to religious leaders than the previous title.
+
+The title bishop is used for both the general (international leader) and the district (state) leaders. The title is sometimes used in conjunction with the previous, thus becoming general (district) superintendent/bishop.
+
+According to the Seventh-day Adventist understanding of the doctrine of the Church:
+
+"The "elders" (Greek, presbuteros) or "bishops" (episkopos) were the most important officers of the church. The term elder means older one, implying dignity and respect. His position was similar to that of the one who had supervision of the synagogue. The term bishop means "overseer." Paul used these terms interchangeably, equating elders with overseers or bishops (Acts 20:17,28; Titus 1:5, 7).
+
+"Those who held this position supervised the newly formed churches. Elder referred to the status or rank of the office, while bishop denoted the duty or responsibility of the office—"overseer." Since the apostles also called themselves elders (1 Peter 5:1; 2 John 1; 3 John 1), it is apparent that there were both local elders and itinerant elders, or elders at large. But both kinds of elder functioned as shepherds of the congregations.[36]"
+
+The above understanding is part of the basis of Adventist organizational structure. The world wide Seventh-day Adventist church is organized into local districts, conferences or missions, union conferences or union missions, divisions, and finally at the top is the general conference.  At each level (with exception to the local districts), there is an elder who is elected president and a group of elders who serve on the executive committee with the elected president.  Those who have been elected president would in effect be the "bishop" while never actually carrying the title or ordained as such because the term is usually associated with the episcopal style of church governance most often found in Catholic, Anglican, Methodist and some Pentecostal/Charismatic circles.
+
+Some Baptists also have begun taking on the title of bishop.[37]
+In some smaller Protestant denominations and independent churches, the term bishop is used in the same way as pastor, to refer to the leader of the local congregation, and may be male or female. This usage is especially common in African-American churches in the US.
+
+In the Church of Scotland, which has a Presbyterian church structure, the word "bishop" refers to an ordained person, usually a normal parish minister, who has temporary oversight of a trainee minister. In the Presbyterian Church (USA), the term bishop is an expressive name for a Minister of Word and Sacrament who serves a congregation and exercises "the oversight of the flock of Christ."[38] The term is traceable to the 1789 Form of Government of the PC (USA) and the Presbyterian understanding of the pastoral office.[39]
+
+While not considered orthodox Christian, the Ecclesia Gnostica Catholica uses roles and titles derived from Christianity for its clerical hierarchy, including bishops who have much the same authority and responsibilities as in Catholicism.
+
+The Salvation Army does not have bishops but has appointed leaders of geographical areas, known as Divisional Commanders. Larger geographical areas, called Territories, are led by a Territorial Commander, who is the highest-ranking officer in that Territory.
+
+Jehovah's Witnesses do not use the title ‘Bishop’ within their organizational structure, but appoint elders to be overseers (to fulfill the role of oversight) within their congregations.[40][41]
+
+The HKBP of Indonesia, the most prominent Protestant denomination in Indonesia, uses the term Ephorus instead of Bishop.[42]
+
+In the Vietnamese syncretist religion of Caodaism, bishops (giáo sư) comprise the fifth of nine hierarchical levels, and are responsible for spiritual and temporal education as well as record-keeping and ceremonies in their parishes. At any one time there are seventy-two bishops. Their authority is described in Section I of the text Tân Luật (revealed through seances in December 1926). Caodai bishops wear robes and headgear of embroidered silk depicting the Divine Eye and the Eight Trigrams. (The color varies according to branch.) This is the full ceremonial dress; the simple version consists of a seven-layered turban.
+
+Traditionally, a number of items are associated with the office of a bishop, most notably the mitre, crosier, and ecclesiastical ring. Other vestments and insignia vary between Eastern and Western Christianity.
+
+In the Latin Rite of the Catholic Church, the choir dress of a bishop includes the purple cassock with amaranth trim, rochet, purple zucchetto (skull cap), purple biretta, and pectoral cross. The cappa magna may be worn, but only within the bishop's own diocese and on especially solemn occasions.[43] The mitre, zuchetto, and stole are generally worn by bishops when presiding over liturgical functions. For liturgical functions other than the Mass the bishop typically wears the cope. Within his own diocese and when celebrating solemnly elsewhere with the consent of the local ordinary, he also uses the crosier.[43] When celebrating Mass, a bishop, like a priest, wears the chasuble. The Caeremoniale Episcoporum recommends, but does not impose, that in solemn celebrations a bishop should also wear a dalmatic, which can always be white, beneath the chasuble, especially when administering the sacrament of holy orders, blessing an abbot or abbess, and dedicating a church or an altar.[43] The Caeremoniale Episcoporum no longer makes mention of episcopal gloves, episcopal sandals, liturgical stockings (also known as buskins), or the accoutrements that it once prescribed for the bishop's horse. The coat of arms of a Latin Rite Catholic bishop usually displays a galero with a cross and crosier behind the escutcheon; the specifics differ by location and ecclesiastical rank (see Ecclesiastical heraldry).
+
+Anglican bishops generally make use of the mitre, crosier, ecclesiastical ring, purple cassock, purple zucchetto, and pectoral cross. However, the traditional choir dress of Anglican bishops retains its late mediaeval form, and looks quite different from that of their Catholic counterparts; it consists of a long rochet which is worn with a chimere.
+
+In the Eastern Churches (Eastern Orthodox, Eastern Rite Catholic) a bishop will wear the mandyas, panagia (and perhaps an enkolpion), sakkos, omophorion and an Eastern-style mitre. Eastern bishops do not normally wear an episcopal ring; the faithful kiss (or, alternatively, touch their forehead to) the bishop's hand. To seal official documents, he will usually use an inked stamp. An Eastern bishop's coat of arms will normally display an Eastern-style mitre, cross, eastern style crosier and a red and white (or red and gold) mantle. The arms of Oriental Orthodox bishops will display the episcopal insignia (mitre or turban) specific to their own liturgical traditions. Variations occur based upon jurisdiction and national customs.
+
+Eastern Rite Catholic bishops celebrating Divine Liturgy in their proper pontifical vestments
+
+An Anglican bishop with a crosier, wearing a rochet under a red chimere and cuffs, a black tippet, and a pectoral cross
+
+An Episcopal bishop immediately before presiding at the Great Vigil of Easter in the narthex of St. Michael's Episcopal Cathedral in Boise, Idaho.
+
+Catholic bishop dressed for the Sacrifice of the Mass. No Pontifical gloves.

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+
+
+Aluminium (aluminum in American and Canadian English) is a chemical element with the symbol Al and atomic number 13. It is a silvery-white, soft, non-magnetic and ductile metal in the boron group. By mass, aluminium makes up about 8% of the Earth's crust, where it is the third most abundant element (after oxygen and silicon) and also the most abundant metal. Occurrence of aluminium decreases in the Earth's mantle below, however. The chief ore of aluminium is bauxite. Aluminium metal is highly reactive, such that native specimens are rare and limited to extreme reducing environments. Instead, it is found combined in over 270 different minerals.[7]
+
+Aluminium is remarkable for its low density and its ability to resist corrosion through the phenomenon of passivation. Aluminium and its alloys are vital to the aerospace industry[8] and important in transportation and building industries, such as building facades and window frames.[9] The oxides and sulfates are the most useful compounds of aluminium.[8]
+
+Despite its prevalence in the environment, no known form of life uses aluminium salts metabolically, but aluminium is well tolerated by plants and animals.[10] Because of these salts' abundance, the potential for a biological role for them is of continuing interest, and studies continue.
+
+Of aluminium isotopes, only 27Al is stable. This is consistent with aluminium having an odd atomic number.[b] It is the only primordial aluminium isotope, i.e. the only one that has existed on Earth in its current form since the creation of the planet. Nearly all aluminium on Earth is present as this isotope, which makes it a mononuclidic element and means that its standard atomic weight is the same as that of the isotope. The standard atomic weight of aluminium is low in comparison with many other metals,[c] which has consequences for the element's properties (see below). This makes aluminium very useful in nuclear magnetic resonance (NMR), as its single stable isotope has a high NMR sensitivity.[12]
+
+All other isotopes of aluminium are radioactive. The most stable of these is 26Al: while it was present along with stable 27Al in the interstellar medium from which the Solar System formed, having been produced by stellar nucleosynthesis as well, its half-life is only 717,000 years and therefore it could not have survived since the formation of the planet. However, minute traces of 26Al are produced from argon in the atmosphere by spallation caused by cosmic ray protons. The ratio of 26Al to 10Be has been used for radiodating of geological processes over 105 to 106 year time scales, in particular transport, deposition, sediment storage, burial times, and erosion.[13] Most meteorite scientists believe that the energy released by the decay of 26Al was responsible for the melting and differentiation of some asteroids after their formation 4.55 billion years ago.[14]
+
+The remaining isotopes of aluminium, with mass numbers ranging from 22 to 43, all have half-lives well under an hour. Three metastable states are known, all with half-lives under a minute.[11]
+
+An aluminium atom has 13 electrons, arranged in an electron configuration of [Ne] 3s2 3p1,[15] with three electrons beyond a stable noble gas configuration. Accordingly, the combined first three ionization energies of aluminium are far lower than the fourth ionization energy alone.[16] Such an electron configuration is shared with the other well-characterized members of its group, boron, gallium, indium, and thallium; it is also expected for nihonium. Aluminium can relatively easily surrender its three outermost electrons in many chemical reactions (see below). The electronegativity of aluminium is 1.61 (Pauling scale).[17]
+
+A free aluminium atom has a radius of 143 pm.[18] With the three outermost electrons removed, the radius shrinks to 39 pm for a 4-coordinated atom or 53.5 pm for a 6-coordinated atom.[18] At standard temperature and pressure, aluminium atoms (when not affected by atoms of other elements) form a face-centered cubic crystal system bound by metallic bonding provided by atoms' outermost electrons; hence aluminium (at these conditions) is a metal.[19] This crystal system is shared by many other metals, such as lead and copper; the size of a unit cell of aluminium is comparable to that of those other metals.[19] It is however not shared by the other members of its group; boron has ionization energies too high to allow metallization, thallium has a hexagonal close-packed structure, and gallium and indium have unusual structures that are not close-packed like those of aluminium and thallium. Since few electrons are available for metallic bonding, aluminium metal is soft with a low melting point and low electrical resistivity, as is common for post-transition metals.[20]
+
+Aluminium metal has an appearance ranging from silvery white to dull gray, depending on the surface roughness. A fresh film of aluminium serves as a good reflector (approximately 92%) of visible light and an excellent reflector (as much as 98%) of medium and far infrared radiation.
+
+The density of aluminium is 2.70 g/cm3, about 1/3 that of steel, much lower than other commonly encountered metals, making aluminium parts easily identifiable through their lightness.[21] Aluminium's low density compared to most other metals arises from the fact that its nuclei are much lighter, while difference in the unit cell size does not compensate for this difference. The only lighter metals are the metals of groups 1 and 2, which apart from beryllium and magnesium are too reactive for structural use (and beryllium is very toxic).[22] Aluminium is not as strong or stiff as steel, but the low density makes up for this in the aerospace industry and for many other applications where light weight and relatively high strength are crucial.
+
+Pure aluminium is quite soft and lacking in strength. In most applications various aluminium alloys are used instead because of their higher strength and hardness. The yield strength of pure aluminium is 7–11 MPa, while aluminium alloys have yield strengths ranging from 200 MPa to 600 MPa.[23] Aluminium is ductile, with a percent elongation of 50-70%,[24] and malleable allowing it to be easily drawn and extruded. It is also easily machined, and the low melting temperature of 660 °C allows for easy casting.
+
+Aluminium is an excellent thermal and electrical conductor, having 59% the conductivity of copper, both thermal and electrical, while having only 30% of copper's density. Aluminium is capable of superconductivity, with a superconducting critical temperature of 1.2 kelvin and a critical magnetic field of about 100 gauss (10 milliteslas).[25] It is paramagnetic and thus essentially unaffected by static magnetic fields. The high electrical conductivity, however, means that it is strongly affected by changing magnetic field through the induction of eddy currents.
+
+Aluminium combines characteristics of pre- and post-transition metals. Since it has few available electrons for metallic bonding, like its heavier group 13 congeners, it has the characteristic physical properties of a post-transition metal, with longer-than-expected interatomic distances.[20] Furthermore, as Al3+ is a small and highly charged cation, it is strongly polarizing and aluminium compounds tend towards covalency;[26] this behaviour is similar to that of beryllium (Be2+), and the two display an example of a diagonal relationship.[27]
+
+The underlying core under aluminium's valence shell is that of the preceding noble gas, whereas those of its heavier congeners gallium and indium, thallium, and nihonium also include a filled d-subshell and in some cases a filled f-subshell. Hence, the inner electrons of aluminium shield the valence electrons almost completely, unlike those of aluminium's heavier congeners. As such, aluminium is the most electropositive metal in its group. In fact, aluminium's electropositive behavior, high affinity for oxygen, and highly negative standard electrode potential are all more similar to those of scandium, yttrium, lanthanum, and actinium, which like aluminium have three valence electrons outside a noble gas core.[20] Aluminium also bears minor similarities to the metalloid boron in the same group: AlX3 compounds are valence isoelectronic to BX3 compounds (they have the same valence electronic structure), and both behave as Lewis acids and readily form adducts.[28] Additionally, one of the main motifs of boron chemistry is regular icosahedral structures, and aluminium forms an important part of many icosahedral quasicrystal alloys, including the Al–Zn–Mg class.[29]
+
+Aluminium has a high chemical affinity to oxygen, which renders it suitable for use as a reducing agent in the thermite reaction. A fine powder of aluminium metal reacts explosively on contact with liquid oxygen; under normal conditions, however, aluminium forms a thin oxide layer (~ 5 nm at room temperature)[30] that protects the metal from further corrosion by oxygen, water, or dilute acid, a process termed passivation.[26][31] Because of its general resistance to corrosion, aluminium is one of the few metals that retains silvery reflectance in finely powdered form, making it an important component of silver-colored paints.[32] Aluminium is not attacked by oxidizing acids because of its passivation. This allows aluminium to be used to store reagents such as nitric acid, concentrated sulfuric acid, and some organic acids.[10]
+
+In hot concentrated hydrochloric acid, aluminium reacts with water with evolution of hydrogen, and in aqueous sodium hydroxide or potassium hydroxide at room temperature to form aluminates—protective passivation under these conditions is negligible.[33] Aqua regia also dissolves aluminium.[10] Aluminium is corroded by dissolved chlorides, such as common sodium chloride, which is why household plumbing is never made from aluminium.[33] The oxide layer on aluminium is also destroyed by contact with mercury due to amalgamation or with salts of some electropositive metals.[26] As such, the strongest aluminium alloys are less corrosion-resistant due to galvanic reactions with alloyed copper,[23] and aluminium's corrosion resistance is greatly reduced by aqueous salts, particularly in the presence of dissimilar metals.[20]
+
+Aluminium reacts with most nonmetals upon heating, forming compounds such as aluminium nitride (AlN), aluminium sulfide (Al2S3), and the aluminium halides (AlX3). It also forms a wide range of intermetallic compounds involving metals from every group on the periodic table.[26]
+
+The vast majority of compounds, including all aluminium-containing minerals and all commercially significant aluminium compounds, feature aluminium in the oxidation state 3+. The coordination number of such compounds varies, but generally Al3+ is either six- or four-coordinate. Almost all compounds of aluminium(III) are colorless.[26]
+
+In aqueous solution, Al3+ exists as the hexaaqua cation [Al(H2O)6]3+, which has an approximate pKa of 10−5.[12] Such solutions are acidic as this cation can act as a proton donor and progressively hydrolyse until a precipitate of aluminium hydroxide, Al(OH)3, forms. This is useful for clarification of water, as the precipitate nucleates on suspended particles in the water, hence removing them. Increasing the pH even further leads to the hydroxide dissolving again as aluminate, [Al(H2O)2(OH)4]−, is formed.
+
+Aluminium hydroxide forms both salts and aluminates and dissolves in acid and alkali, as well as on fusion with acidic and basic oxides.[26] This behaviour of Al(OH)3 is termed amphoterism, and is characteristic of weakly basic cations that form insoluble hydroxides and whose hydrated species can also donate their protons. One effect of this is that aluminium salts with weak acids are hydrolysed in water to the aquated hydroxide and the corresponding nonmetal hydride: for example, aluminium sulfide yields hydrogen sulfide. However, some salts like aluminium carbonate exist in aqueous solution but are unstable as such; and only incomplete hydrolysis takes place for salts with strong acids, such as the halides, nitrate, and sulfate. For similar reasons, anhydrous aluminium salts cannot be made by heating their "hydrates": hydrated aluminium chloride is in fact not AlCl3·6H2O but [Al(H2O)6]Cl3, and the Al–O bonds are so strong that heating is not sufficient to break them and form Al–Cl bonds instead:[26]
+
+All four trihalides are well known. Unlike the structures of the three heavier trihalides, aluminium fluoride (AlF3) features six-coordinate aluminium, which explains its involatility and insolubility as well as high heat of formation. Each aluminium atom is surrounded by six fluorine atoms in a distorted octahedral arrangement, with each fluorine atom being shared between the corners of two octahedra. Such {AlF6} units also exist in complex fluorides such as cryolite, Na3AlF6.[d] AlF3 melts at 1,290 °C (2,354 °F) and is made by reaction of aluminium oxide with hydrogen fluoride gas at 700 °C (1,292 °F).[35]
+
+With heavier halides, the coordination numbers are lower. The other trihalides are dimeric or polymeric with tetrahedral four-coordinate aluminium centers. Aluminium trichloride (AlCl3) has a layered polymeric structure below its melting point of 192.4 °C (378 °F) but transforms on melting to Al2Cl6 dimers. At higher temperatures those increasingly dissociate into trigonal planar AlCl3 monomers similar to the structure of BCl3. Aluminium tribromide and aluminium triiodide form Al2X6 dimers in all three phases and hence do not show such significant changes of properties upon phase change.[35] These materials are prepared by treating aluminium metal with the halogen. The aluminium trihalides form many addition compounds or complexes; their Lewis acidic nature makes them useful as catalysts for the Friedel–Crafts reactions. Aluminium trichloride has major industrial uses involving this reaction, such as in the manufacture of anthraquinones and styrene; it is also often used as the precursor for many other aluminium compounds and as a reagent for converting nonmetal fluorides into the corresponding chlorides (a transhalogenation reaction).[35]
+
+Aluminium forms one stable oxide with the chemical formula Al2O3, commonly called alumina.[36] It can be found in nature in the mineral corundum, α-alumina;[37] there is also a γ-alumina phase.[12] Its crystalline form, corundum, is very hard (Mohs hardness 9), has a high melting point of 2,045 °C (3,713 °F), has very low volatility, is chemically inert, and a good electrical insulator, it is often used in abrasives (such as toothpaste), as a refractory material, and in ceramics, as well as being the starting material for the electrolytic production of aluminium metal. Sapphire and ruby are impure corundum contaminated with trace amounts of other metals.[12] The two main oxide-hydroxides, AlO(OH), are boehmite and diaspore. There are three main trihydroxides: bayerite, gibbsite, and nordstrandite, which differ in their crystalline structure (polymorphs). Many other intermediate and related structures are also known.[12] Most are produced from ores by a variety of wet processes using acid and base. Heating the hydroxides leads to formation of corundum. These materials are of central importance to the production of aluminium and are themselves extremely useful. Some mixed oxide phases are also very useful, such as spinel (MgAl2O4), Na-β-alumina (NaAl11O17), and tricalcium aluminate (Ca3Al2O6, an important mineral phase in Portland cement).[12]
+
+The only stable chalcogenides under normal conditions are aluminium sulfide (Al2S3), selenide (Al2Se3), and telluride (Al2Te3). All three are prepared by direct reaction of their elements at about 1,000 °C (1,832 °F) and quickly hydrolyse completely in water to yield aluminium hydroxide and the respective hydrogen chalcogenide. As aluminium is a small atom relative to these chalcogens, these have four-coordinate tetrahedral aluminium with various polymorphs having structures related to wurtzite, with two-thirds of the possible metal sites occupied either in an orderly (α) or random (β) fashion; the sulfide also has a γ form related to γ-alumina, and an unusual high-temperature hexagonal form where half the aluminium atoms have tetrahedral four-coordination and the other half have trigonal bipyramidal five-coordination.[38]
+
+Four pnictides – aluminium nitride (AlN), aluminium phosphide (AlP), aluminium arsenide (AlAs), and aluminium antimonide (AlSb) – are known. They are all III-V semiconductors isoelectronic to silicon and germanium, all of which but AlN have the zinc blende structure. All four can be made by high-temperature (and possibly high-pressure) direct reaction of their component elements.[38]
+
+Although the great majority of aluminium compounds feature Al3+ centers, compounds with lower oxidation states are known and are sometimes of significance as precursors to the Al3+ species.
+
+AlF, AlCl, AlBr, and AlI exist in the gaseous phase when the respective trihalide is heated with aluminium, and at cryogenic temperatures. Their instability in the condensed phase is due to their ready disproportionation to aluminium and the respective trihalide: the reverse reaction is favored at high temperature (although even then they are still short-lived), explaining why AlF3 is more volatile when heated in the presence of aluminium metal, as is aluminium metal when heated in the presence of AlCl3.[35] A stable derivative of aluminium monoiodide is the cyclic adduct formed with triethylamine, Al4I4(NEt3)4. Also of theoretical interest but only of fleeting existence are Al2O and Al2S. Al2O is made by heating the normal oxide, Al2O3, with silicon at 1,800 °C (3,272 °F) in a vacuum. Such materials quickly disproportionate to the starting materials.[39]
+
+Very simple Al(II) compounds are invoked or observed in the reactions of Al metal with oxidants. For example, aluminium monoxide, AlO, has been detected in the gas phase after explosion[40] and in stellar absorption spectra.[41] More thoroughly investigated are compounds of the formula R4Al2 which contain an Al–Al bond and where R is a large organic ligand.[42]
+
+A variety of compounds of empirical formula AlR3 and AlR1.5Cl1.5 exist.[43] The aluminium trialkyls and triaryls are reactive, volatile, and colorless liquids or low-melting solids. They catch fire spontaneously in air and react with water, thus necessitating precautions when handling them. They often form dimers, unlike their boron analogues, but this tendency diminishes for branched-chain alkyls (e.g. Pri, Bui, Me3CCH2); for example, triisobutylaluminium exists as an equilibrium mixture of the monomer and dimer.[44][45] These dimers, such as trimethylaluminium (Al2Me6), usually feature tetrahedral Al centers formed by dimerization with some alkyl group bridging between both aluminium atoms. They are hard acids and react readily with ligands, forming adducts. In industry, they are mostly used in alkene insertion reactions, as discovered by Karl Ziegler, most importantly in "growth reactions" that form long-chain unbranched primary alkenes and alcohols, and in the low-pressure polymerization of ethene and propene. There are also some heterocyclic and cluster organoaluminium compounds involving Al–N bonds.[44]
+
+The industrially most important aluminium hydride is lithium aluminium hydride (LiAlH4), which is used in as a reducing agent in organic chemistry. It can be produced from lithium hydride and aluminium trichloride.[46] The simplest hydride, aluminium hydride or alane, is not as important. It is a polymer with the formula (AlH3)n, in contrast to the corresponding boron hydride that is a dimer with the formula (BH3)2.[46]
+
+Aluminium's per-particle abundance in the Solar System is 3.15 ppm (parts per million).[47][e] It is the twelfth most abundant of all elements and third most abundant among the elements that have odd atomic numbers, after hydrogen and nitrogen.[47] The only stable isotope of aluminium, 27Al, is the eighteenth most abundant nucleus in the Universe. It is created almost entirely after fusion of carbon in massive stars that will later become Type II supernovae: this fusion creates 26Mg, which, upon capturing free protons and neutrons becomes aluminium. Some smaller quantities of 27Al are created in hydrogen burning shells of evolved stars, where 26Mg can capture free protons.[48] Essentially all aluminium now in existence is 27Al; 26Al was present in the early Solar System but is currently extinct. However, the trace quantities of 26Al that do exist are the most common gamma ray emitter in the interstellar gas.[48]
+
+Overall, the Earth is about 1.59% aluminium by mass (seventh in abundance by mass).[49] Aluminium occurs in greater proportion in the Earth than in the Universe because aluminium easily forms the oxide and becomes bound into rocks and aluminium stays in the Earth's crust while less reactive metals sink to the core.[48] In the Earth's crust, aluminium is the most abundant (8.23% by mass[24]) metallic element and the third most abundant of all elements (after oxygen and silicon).[50] A large number of silicates in the Earth's crust contain aluminium.[51] In contrast, the Earth's mantle is only 2.38% aluminium by mass.[52] Aluminium also occurs in seawater at a concentration of 2 μg/kg.[24]
+
+Because of its strong affinity for oxygen, aluminium is almost never found in the elemental state; instead it is found in oxides or silicates. Feldspars, the most common group of minerals in the Earth's crust, are aluminosilicates. Aluminium also occurs in the minerals beryl, cryolite, garnet, spinel, and turquoise.[53] Impurities in Al2O3, such as chromium and iron, yield the gemstones ruby and sapphire, respectively.[54] Native aluminium metal can only be found as a minor phase in low oxygen fugacity environments, such as the interiors of certain volcanoes.[55] Native aluminium has been reported in cold seeps in the northeastern continental slope of the South China Sea. It is possible that these deposits resulted from bacterial reduction of tetrahydroxoaluminate Al(OH)4−.[56]
+
+Although aluminium is a common and widespread element, not all aluminium minerals are economically viable sources of the metal. Almost all metallic aluminium is produced from the ore bauxite (AlOx(OH)3–2x). Bauxite occurs as a weathering product of low iron and silica bedrock in tropical climatic conditions.[57] In 2017, most bauxite was mined in Australia, China, Guinea, and India.[58]
+
+The history of aluminium has been shaped by usage of alum. The first written record of alum, made by Greek historian Herodotus, dates back to the 5th century BCE.[59] The ancients are known to have used alum as a dyeing mordant and for city defense.[59] After the Crusades, alum, an indispensable good in the European fabric industry,[60] was a subject of international commerce;[61] it was imported to Europe from the eastern Mediterranean until the mid-15th century.[62]
+
+The nature of alum remained unknown. Around 1530, Swiss physician Paracelsus suggested alum was a salt of an earth of alum.[63] In 1595, German doctor and chemist Andreas Libavius experimentally confirmed this.[64] In 1722, German chemist Friedrich Hoffmann announced his belief that the base of alum was a distinct earth.[65] In 1754, German chemist Andreas Sigismund Marggraf synthesized alumina by boiling clay in sulfuric acid and subsequently adding potash.[65]
+
+Attempts to produce aluminium metal date back to 1760.[66] The first successful attempt, however, was completed in 1824 by Danish physicist and chemist Hans Christian Ørsted. He reacted anhydrous aluminium chloride with potassium amalgam, yielding a lump of metal looking similar to tin.[67][68][69] He presented his results and demonstrated a sample of the new metal in 1825.[70][71] In 1827, German chemist Friedrich Wöhler repeated Ørsted's experiments but did not identify any aluminium.[72] (The reason for this inconsistency was only discovered in 1921.)[73] He conducted a similar experiment in the same year by mixing anhydrous aluminium chloride with potassium and produced a powder of aluminium.[69] In 1845, he was able to produce small pieces of the metal and described some physical properties of this metal.[73] For many years thereafter, Wöhler was credited as the discoverer of aluminium.[74]
+
+As Wöhler's method could not yield great quantities of aluminium, the metal remained rare; its cost exceeded that of gold.[72] The first industrial production of aluminium was established in 1856 by French chemist Henri Etienne Sainte-Claire Deville and companions.[75] Deville had discovered that aluminium trichloride could be reduced by sodium, which was more convenient and less expensive than potassium, which Wöhler had used.[76] Even then, aluminium was still not of great purity and produced aluminium differed in properties by sample.[77]
+
+The first industrial large-scale production method was independently developed in 1886 by French engineer Paul Héroult and American engineer Charles Martin Hall; it is now known as the Hall–Héroult process.[78] The Hall–Héroult process converts alumina into the metal. Austrian chemist Carl Joseph Bayer discovered a way of purifying bauxite to yield alumina, now known as the Bayer process, in 1889.[79] Modern production of the aluminium metal is based on the Bayer and Hall–Héroult processes.[80]
+
+Prices of aluminium dropped and aluminium became widely used in jewelry, everyday items, eyeglass frames, optical instruments, tableware, and foil in the 1890s and early 20th century. Aluminium's ability to form hard yet light alloys with other metals provided the metal many uses at the time.[81] During World War I, major governments demanded large shipments of aluminium for light strong airframes.[82]
+
+By the mid-20th century, aluminium had become a part of everyday life and an essential component of housewares.[83] During the mid-20th century, aluminium emerged as a civil engineering material, with building applications in both basic construction and interior finish work,[84] and increasingly being used in military engineering, for both airplanes and land armor vehicle engines.[85] Earth's first artificial satellite, launched in 1957, consisted of two separate aluminium semi-spheres joined together and all subsequent space vehicles have used aluminium to some extent.[80] The aluminium can was invented in 1956 and employed as a storage for drinks in 1958.[86]
+
+Throughout the 20th century, the production of aluminium rose rapidly: while the world production of aluminium in 1900 was 6,800 metric tons, the annual production first exceeded 100,000 metric tons in 1916; 1,000,000 tons in 1941; 10,000,000 tons in 1971.[87] In the 1970s, the increased demand for aluminium made it an exchange commodity; it entered the London Metal Exchange, the oldest industrial metal exchange in the world, in 1978.[80] The output continued to grow: the annual production of aluminium exceeded 50,000,000 metric tons in 2013.[87]
+
+The real price for aluminium declined from $14,000 per metric ton in 1900 to $2,340 in 1948 (in 1998 United States dollars).[87] Extraction and processing costs were lowered over technological progress and the scale of the economies. However, the need to exploit lower-grade poorer quality deposits and the use of fast increasing input costs (above all, energy) increased the net cost of aluminium;[88] the real price began to grow in the 1970s with the rise of energy cost.[89] Production moved from the industrialized countries to countries where production was cheaper.[90] Production costs in the late 20th century changed because of advances in technology, lower energy prices, exchange rates of the United States dollar, and alumina prices.[91] The BRIC countries' combined share in primary production and primary consumption grew substantially in the first decade of the 21st century.[92] China is accumulating an especially large share of world's production thanks to abundance of resources, cheap energy, and governmental stimuli;[93] it also increased its consumption share from 2% in 1972 to 40% in 2010.[94] In the United States, Western Europe, and Japan, most aluminium was consumed in transportation, engineering, construction, and packaging.[95]
+
+Aluminium is named after alumina, or aluminium oxide in modern nomenclature. The word "alumina" comes from "alum", the mineral from which it was collected. The word "alum" comes from alumen, a Latin word meaning "bitter salt".[96] The word alumen stems from the Proto-Indo-European root *alu- meaning "bitter" or "beer".[97]
+
+British chemist Humphry Davy, who performed a number of experiments aimed to isolate the metal, is credited as the person who named the element. In 1808, he suggested the metal be named alumium in an article on his electrochemical research which was published in Philosophical Transactions of the Royal Society.[98] This suggestion was criticized by contemporary chemists from France, Germany, and Sweden, who insisted the metal should be named for the oxide, alumina, from which it would be isolated.[99] In 1812, Davy published a chemistry textbook in which he settled on the name aluminum, thus producing the modern name.[100] However, its spelling and pronunciation varies: aluminum is in use in the United States and Canada while aluminium is in use elsewhere.[101]
+
+The -ium suffix followed the precedent set in other newly discovered elements of the time: potassium, sodium, magnesium, calcium, and strontium (all of which Davy isolated himself). Nevertheless, element names ending in -um were known at the time, for example, platinum (known to Europeans since the 16th century), molybdenum (discovered in 1778), and tantalum (discovered in 1802). The -um suffix is consistent with the universal spelling alumina for the oxide (as opposed to aluminia); compare to lanthana, the oxide of lanthanum, and magnesia, ceria, and thoria, the oxides of magnesium, cerium, and thorium, respectively.
+
+In 1812, British scientist Thomas Young[102] wrote an anonymous review of Davy's book, in which he objected to aluminum and proposed the name aluminium: "for so we shall take the liberty of writing the word, in preference to aluminum, which has a less classical sound."[103] This name did catch on: while the -um spelling was occasionally used in Britain, the American scientific language used -ium from the start.[104] Most scientists used -ium throughout the world in the 19th century;[105] it still remains the standard in most other languages.[101] In 1828, American lexicographer Noah Webster used exclusively the aluminum spelling in his American Dictionary of the English Language.[106] In the 1830s, the -um spelling started to gain usage in the United States; by the 1860s, it had become the more common spelling there outside science.[104] In 1892, Hall used the -um spelling in his advertising handbill for his new electrolytic method of producing the metal, despite his constant use of the -ium spelling in all the patents he filed between 1886 and 1903. It was subsequently suggested this was a typo rather than intended.[101] By 1890, both spellings had been common in the U.S. overall, the -ium spelling being slightly more common; by 1895, the situation had reversed; by 1900, aluminum had become twice as common as aluminium; during the following decade, the -um spelling dominated American usage.[107] In 1925, the American Chemical Society adopted this spelling.[107]
+
+The International Union of Pure and Applied Chemistry (IUPAC) adopted aluminium as the standard international name for the element in 1990.[108] In 1993, they recognized aluminum as an acceptable variant;[108] the most recent 2005 edition of the IUPAC nomenclature of inorganic chemistry acknowledges this spelling as well.[109] IUPAC official publications use the -ium spelling as primary but list both where appropriate.[f]
+
+Aluminium production is highly energy-consuming, and so the producers tend to locate smelters in places where electric power is both plentiful and inexpensive.[112] As of 2012, the world's largest smelters of aluminium are located in China, Russia, Bahrain, United Arab Emirates, and South Africa.[113]
+
+In 2016, China was the top producer of aluminium with a world share of fifty-five percent; the next largest producing countries were Russia, Canada, India, and the United Arab Emirates.[111]
+
+According to the International Resource Panel's Metal Stocks in Society report, the global per capita stock of aluminium in use in society (i.e. in cars, buildings, electronics, etc.) is 80 kg (180 lb). Much of this is in more-developed countries (350–500 kg (770–1,100 lb) per capita) rather than less-developed countries (35 kg (77 lb) per capita).[114]
+
+Bauxite is converted to aluminium oxide by the Bayer process. Bauxite is blended for uniform composition and then is ground. The resulting slurry is mixed with a hot solution of sodium hydroxide; the mixture is then treated in a digester vessel at a pressure well above atmospheric, dissolving the aluminium hydroxide in bauxite while converting impurities into relatively insoluble compounds:[115]
+
+After this reaction, the slurry is at a temperature above its atmospheric boiling point. It is cooled by removing steam as pressure is reduced. The bauxite residue is separated from the solution and discarded. The solution, free of solids, is seeded with small crystals of aluminium hydroxide; this causes decomposition of the [Al(OH)4]− ions to aluminium hydroxide. After about half of aluminium has precipitated, the mixture is sent to classifiers. Small crystals of aluminium hydroxide are collected to serve as seeding agents; coarse particles are converted to aluminium oxide by heating; excess solution is removed by evaporation, (if needed) purified, and recycled.[115]
+
+The conversion of alumina to aluminium metal is achieved by the Hall–Héroult process. In this energy-intensive process, a solution of alumina in a molten (950 and 980 °C (1,740 and 1,800 °F)) mixture of cryolite (Na3AlF6) with calcium fluoride is electrolyzed to produce metallic aluminium. The liquid aluminium metal sinks to the bottom of the solution and is tapped off, and usually cast into large blocks called aluminium billets for further processing.[10]
+
+Anodes of the electrolysis cell are made of carbon—the most resistant material against fluoride corrosion—and either bake at the process or are prebaked. The former, also called Söderberg anodes, are less power-efficient and fumes released during baking are costly to collect, which is why they are being replaced by prebaked anodes even though they save the power, energy, and labor to prebake the cathodes. Carbon for anodes should be preferably pure so that neither aluminium nor the electrolyte is contaminated with ash. Despite carbon's resistivity against corrosion, it is still consumed at a rate of 0.4–0.5 kg per each kilogram of produced aluminium. Cathodes are made of anthracite; high purity for them is not required because impurities leach only very slowly. Cathode is consumed at a rate of 0.02–0.04 kg per each kilogram of produced aluminium. A cell is usually a terminated after 2–6 years following a failure of the cathode.[10]
+
+The Hall–Heroult process produces aluminium with a purity of above 99%. Further purification can be done by the Hoopes process. This process involves the electrolysis of molten aluminium with a sodium, barium, and aluminium fluoride electrolyte. The resulting aluminium has a purity of 99.99%.[10][116]
+
+Electric power represents about 20 to 40% of the cost of producing aluminium, depending on the location of the smelter. Aluminium production consumes roughly 5% of electricity generated in the United States.[108] Because of this, alternatives to the Hall–Héroult process have been researched, but none has turned out to be economically feasible.[10]
+
+Recovery of the metal through recycling has become an important task of the aluminium industry. Recycling was a low-profile activity until the late 1960s, when the growing use of aluminium beverage cans brought it to public awareness.[117] Recycling involves melting the scrap, a process that requires only 5% of the energy used to produce aluminium from ore, though a significant part (up to 15% of the input material) is lost as dross (ash-like oxide).[118] An aluminium stack melter produces significantly less dross, with values reported below 1%.[119]
+
+White dross from primary aluminium production and from secondary recycling operations still contains useful quantities of aluminium that can be extracted industrially. The process produces aluminium billets, together with a highly complex waste material. This waste is difficult to manage. It reacts with water, releasing a mixture of gases (including, among others, hydrogen, acetylene, and ammonia), which spontaneously ignites on contact with air;[120] contact with damp air results in the release of copious quantities of ammonia gas. Despite these difficulties, the waste is used as a filler in asphalt and concrete.[121]
+
+Aluminium is the most widely used non-ferrous metal.[122] The global production of aluminium in 2016 was 58.8 million metric tons. It exceeded that of any other metal except iron (1,231 million metric tons).[111]
+
+Aluminium is almost always alloyed, which markedly improves its mechanical properties, especially when tempered. For example, the common aluminium foils and beverage cans are alloys of 92% to 99% aluminium.[123] The main alloying agents are copper, zinc, magnesium, manganese, and silicon (e.g., duralumin) with the levels of other metals in a few percent by weight.[124]
+
+The major uses for aluminium metal are in:[125]
+
+The great majority (about 90%) of aluminium oxide is converted to metallic aluminium.[115] Being a very hard material (Mohs hardness 9),[126] alumina is widely used as an abrasive;[127] being extraordinarily chemically inert, it is useful in highly reactive environments such as high pressure sodium lamps.[128] Aluminium oxide is commonly used as a catalyst for industrial processes;[115] e.g. the Claus process to convert hydrogen sulfide to sulfur in refineries and to alkylate amines.[129][130] Many industrial catalysts are supported by alumina, meaning that the expensive catalyst material is dispersed over a surface of the inert alumina.[131] Another principal use is as a drying agent or absorbent.[115][132]
+
+Several sulfates of aluminium have industrial and commercial application. Aluminium sulfate (in its hydrate form) is produced on the annual scale of several millions of metric tons.[133] About two-thirds is consumed in water treatment.[133] The next major application is in the manufacture of paper.[133] It is also used as a mordant in dyeing, in pickling seeds, deodorizing of mineral oils, in leather tanning, and in production of other aluminium compounds.[133] Two kinds of alum, ammonium alum and potassium alum, were formerly used as mordants and in leather tanning, but their use has significantly declined following availability of high-purity aluminium sulfate.[133] Anhydrous aluminium chloride is used as a catalyst in chemical and petrochemical industries, the dyeing industry, and in synthesis of various inorganic and organic compounds.[133] Aluminium hydroxychlorides are used in purifying water, in the paper industry, and as antiperspirants.[133] Sodium aluminate is used in treating water and as an accelerator of solidification of cement.[133]
+
+Many aluminium compounds have niche applications, for example:
+
+Despite its widespread occurrence in the Earth's crust, aluminium has no known function in biology.[10] At pH 6–9 (relevant for most natural waters), aluminium precipitates out of water as the hydroxide and is hence not available; most elements behaving this way have no biological role or are toxic.[145] Aluminium salts are remarkably nontoxic, aluminium sulfate having an LD50 of 6207 mg/kg (oral, mouse), which corresponds to 500 grams for an 80 kg (180 lb) person.[10]
+
+In most people, aluminium is not as toxic as heavy metals. Aluminium is classified as a non-carcinogen by the United States Department of Health and Human Services.[146] There is little evidence that normal exposure to aluminium presents a risk to healthy adult,[147] and there is evidence of no toxicity if it is consumed in amounts not greater than 40 mg/day per kg of body mass.[146] Most aluminium consumed will leave the body in feces; most of the small part of it that enters the bloodstream, will be excreted via urine.[148]
+
+Aluminium, although rarely, can cause vitamin D-resistant osteomalacia, erythropoietin-resistant microcytic anemia, and central nervous system alterations. People with kidney insufficiency are especially at a risk.[146] Chronic ingestion of hydrated aluminium silicates (for excess gastric acidity control) may result in aluminium binding to intestinal contents and increased elimination of other metals, such as iron or zinc; sufficiently high doses (>50 g/day) can cause anemia.[146]
+
+During the 1988 Camelford water pollution incident people in Camelford had their drinking water contaminated with aluminium sulfate for several weeks. A final report into the incident in 2013 concluded it was unlikely that this had caused long-term health problems.[149]
+
+Aluminium has been suspected of being a possible cause of Alzheimer's disease,[150] but research into this for over 40 years has found, as of 2018[update], no good evidence of causal effect.[151][152]
+
+Aluminium increases estrogen-related gene expression in human breast cancer cells cultured in the laboratory.[153] In very high doses, aluminium is associated with altered function of the blood–brain barrier.[154] A small percentage of people[155] have contact allergies to aluminium and experience itchy red rashes, headache, muscle pain, joint pain, poor memory, insomnia, depression, asthma, irritable bowel syndrome, or other symptoms upon contact with products containing aluminium.[156]
+
+Exposure to powdered aluminium or aluminium welding fumes can cause pulmonary fibrosis.[157] Fine aluminium powder can ignite or explode, posing another workplace hazard.[158][159]
+
+Food is the main source of aluminium. Drinking water contains more aluminium than solid food;[146] however, aluminium in food may be absorbed more than aluminium from water.[160] Major sources of human oral exposure to aluminium include food (due to its use in food additives, food and beverage packaging, and cooking utensils), drinking water (due to its use in municipal water treatment), and aluminium-containing medications (particularly antacid/antiulcer and buffered aspirin formulations).[161] Dietary exposure in Europeans averages to 0.2–1.5 mg/kg/week but can be as high as 2.3 mg/kg/week.[146] Higher exposure levels of aluminium are mostly limited to miners, aluminium production workers, and dialysis patients.[162]
+
+Consumption of antacids, antiperspirants, vaccines, and cosmetics provide possible routes of exposure.[163] Consumption of acidic foods or liquids with aluminium enhances aluminium absorption,[164] and maltol has been shown to increase the accumulation of aluminium in nerve and bone tissues.[165]
+
+In case of suspected sudden intake of a large amount of aluminium, the only treatment is deferoxamine mesylate which may be given to help eliminate aluminium from the body by chelation.[166][167] However, this should be applied with caution as this reduces not only aluminium body levels, but also those of other metals such as copper or iron.[166]
+
+High levels of aluminium occur near mining sites; small amounts of aluminium are released to the environment at the coal-fired power plants or incinerators.[168] Aluminium in the air is washed out by the rain or normally settles down but small particles of aluminium remain in the air for a long time.[168]
+
+Acidic precipitation is the main natural factor to mobilize aluminium from natural sources[146] and the main reason for the environmental effects of aluminium;[169] however, the main factor of presence of aluminium in salt and freshwater are the industrial processes that also release aluminium into air.[146]
+
+In water, aluminium acts as a toxiс agent on gill-breathing animals such as fish by causing loss of plasma- and hemolymph ions leading to osmoregulatory failure.[169] Organic complexes of aluminium may be easily absorbed and interfere with metabolism in mammals and birds, even though this rarely happens in practice.[169]
+
+Aluminium is primary among the factors that reduce plant growth on acidic soils. Although it is generally harmless to plant growth in pH-neutral soils, in acid soils the concentration of toxic Al3+ cations increases and disturbs root growth and function.[170][171][172][173] Wheat has developed a tolerance to aluminium, releasing organic compounds that bind to harmful aluminium cations. Sorghum is believed to have the same tolerance mechanism.[174]
+
+Aluminium production possesses its own challenges to the environment on each step of the production process. The major challenge is the greenhouse gas emissions.[162] These gases result from electrical consumption of the smelters and the byproducts of processing. The most potent of these gases are perfluorocarbons from the smelting process.[162] Released sulfur dioxide is one of the primary precursors of acid rain.[162]
+
+A Spanish scientific report from 2001 claimed that the fungus Geotrichum candidum consumes the aluminium in compact discs.[175][176] Other reports all refer back to that report and there is no supporting original research. Better documented, the bacterium Pseudomonas aeruginosa and the fungus Cladosporium resinae are commonly detected in aircraft fuel tanks that use kerosene-based fuels (not avgas), and laboratory cultures can degrade aluminium.[177] However, these life forms do not directly attack or consume the aluminium; rather, the metal is corroded by microbe waste products.[178]

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+Homer (/ˈhoʊmər/; Ancient Greek: Ὅμηρος Greek pronunciation: [hómɛːros], Hómēros) is the presumed author of the Iliad and the Odyssey, two epic poems that are the central works of ancient Greek literature. The Iliad is set during the Trojan War, the ten-year siege of the city of Troy by a coalition of Greek kingdoms. It focuses on a quarrel between King Agamemnon and the warrior Achilles lasting a few weeks during the last year of the war. The Odyssey focuses on the ten-year journey home of Odysseus, king of Ithaca, after the fall of Troy. Many accounts of Homer's life circulated in classical antiquity, the most widespread being that he was a blind bard from Ionia, a region of central coastal Anatolia in present-day Turkey. Modern scholars consider these accounts legendary.[2][3][4]
+
+The Homeric Question – concerning by whom, when, where and under what circumstances the Iliad and Odyssey were composed – continues to be debated. Broadly speaking, modern scholarly opinion falls into two groups. One holds that most of the Iliad and (according to some) the Odyssey are the works of a single poet of genius. The other considers the Homeric poems to be the result of a process of working and reworking by many contributors, and that "Homer" is best seen as a label for an entire tradition.[4] It is generally accepted that the poems were composed at some point around the late eighth or early seventh century BC.[5]
+
+The poems are in Homeric Greek, also known as Epic Greek, a literary language which shows a mixture of features of the Ionic and Aeolic dialects from different centuries; the predominant influence is Eastern Ionic.[6][7] Most researchers believe that the poems were originally transmitted orally.[8] From antiquity until the present day, the influence of Homeric epic on Western civilization has been great, inspiring many of its most famous works of literature, music, art and film.[9] The Homeric epics were the greatest influence on ancient Greek culture and education; to Plato, Homer was simply the one who "has taught Greece"  – ten Hellada pepaideuken.[10][11]
+
+Today only the Iliad and Odyssey are associated with the name 'Homer'. In antiquity, a very large number of other works were sometimes attributed to him, including the Homeric Hymns, the Contest of Homer and Hesiod, the Little Iliad, the Nostoi, the Thebaid, the Cypria, the Epigoni, the comic mini-epic Batrachomyomachia ("The Frog-Mouse War"), the Margites, the Capture of Oechalia, and the Phocais. These claims are not considered authentic today and were by no means universally accepted in the ancient world. As with the multitude of legends surrounding Homer's life, they indicate little more than the centrality of Homer to ancient Greek culture.[12][13][14]
+
+Many traditions circulated in the ancient world concerning Homer, most of which are lost. Modern scholarly consensus is that they have no value as history. Some claims were established early and repeated often. They include that Homer was blind (taking as self-referential a passage describing the blind bard Demodocus[15][16]), that he was born in Chios, that he was the son of the river Meles and the nymph Critheïs, that he was a wandering bard, that he composed a varying list of other works (the "Homerica"), that he died either in Ios or after failing to solve a riddle set by fishermen, and various explanations for the name "Homer". The two best known ancient biographies of Homer are the Life of Homer by the Pseudo-Herodotus and the Contest of Homer and Hesiod.[17][18]
+
+The study of Homer is one of the oldest topics in scholarship, dating back to antiquity.[19][20][21] Nonetheless, the aims of Homeric studies have changed over the course of the millennia.[19] The earliest preserved comments on Homer concern his treatment of the gods, which hostile critics such as the poet Xenophanes of Colophon denounced as immoral.[21] The allegorist Theagenes of Rhegium is said to have defended Homer by arguing that the Homeric poems are allegories.[21] The Iliad and the Odyssey were widely used as school texts in ancient Greek and Hellenistic cultures.[19][21][22] They were the first literary works taught to all students.[22] The Iliad, particularly its first few books, was far more intently studied than the Odyssey during the Hellenistic and Roman periods.[22]
+
+As a result of the poems' prominence in classical Greek education, extensive commentaries on them developed to explain parts of the poems that were culturally or linguistically difficult.[19][21] During the Hellenistic and Roman periods, many interpreters, especially the Stoics, who believed that Homeric poems conveyed Stoic doctrines, regarded them as allegories, containing hidden wisdom.[21] Perhaps partially because of the Homeric poems' extensive use in education, many authors believed that Homer's original purpose had been to educate.[21] Homer's wisdom became so widely praised that he began to acquire the image of almost a prototypical philosopher.[21] Byzantine scholars such as Eustathius of Thessalonica and John Tzetzes produced commentaries, extensions and scholia to Homer, especially in the twelfth century.[23][21] Eustathius's commentary on the Iliad alone is massive, sprawling over nearly 4,000 oversized pages in a twenty-first century printed version and his commentary on the Odyssey an additional nearly 2,000.[21]
+
+In 1488, the Greek scholar Demetrios Chalkokondyles published the editio princeps of the Homeric poems.[21] The earliest modern Homeric scholars started with the same basic approaches towards the Homeric poems as scholars in antiquity.[21][20][19] The allegorical interpretation of the Homeric poems that had been so prevalent in antiquity returned to become the prevailing view of the Renaissance.[21] Renaissance humanists praised Homer as the archetypically wise poet, whose writings contain hidden wisdom, disguised through allegory.[21] In western Europe during the Renaissance, Virgil was more widely read than Homer and Homer was often seen through a Virgilian lens.[24]
+
+In 1664, contradicting the widespread praise of Homer as the epitome of wisdom, François Hédelin, abbé d'Aubignac wrote a scathing attack on the Homeric poems, declaring that they were incoherent, immoral, tasteless, and without style, that Homer never existed, and that the poems were hastily cobbled together by incompetent editors from unrelated oral songs.[20] Fifty years later, the English scholar Richard Bentley concluded that Homer did exist, but that he was an obscure, prehistoric oral poet whose compositions bear little relation to the Iliad and the Odyssey as they have been passed down.[20] According to Bentley, Homer "wrote a Sequel of Songs and Rhapsodies, to be sung by himself for small Earnings and good Cheer at Festivals and other Days of Merriment; the Ilias he wrote for men, and the Odysseis for the other Sex. These loose songs were not collected together in the Form of an epic Poem till Pisistratus' time, about 500 Years after."[20]
+
+Friedrich August Wolf's Prolegomena ad Homerum, published in 1795, argued that much of the material later incorporated into the Iliad and the Odyssey was originally composed in the tenth century BC in the form of short, separate oral songs,[25][26][20] which passed through oral tradition for roughly four hundred years before being assembled into prototypical versions of the Iliad and the Odyssey in the sixth century BC by literate authors.[25][26][20] After being written down, Wolf maintained that the two poems were extensively edited, modernized, and eventually shaped into their present state as artistic unities.[25][26][20] Wolf and the "Analyst" school, which led the field in the nineteenth century, sought to recover the original, authentic poems which were thought to be concealed by later excrescences.[25][26][20][27]
+
+Within the Analyst school were two camps: proponents of the "lay theory," which held that the Iliad and the Odyssey were put together from a large number of short, independent songs,[20] and proponents of the "nucleus theory", which held that Homer had originally composed shorter versions of the Iliad and the Odyssey, which later poets expanded and revised.[20] A small group of scholars opposed to the Analysts, dubbed "Unitarians", saw the later additions as superior, the work of a single inspired poet.[25][26][20] By around 1830, the central preoccupations of Homeric scholars, dealing with whether or not "Homer" actually existed, when and how the Homeric poems originated, how they were transmitted, when and how they were finally written down, and their overall unity, had been dubbed "the Homeric Question".[20]
+
+Following World War I, the Analyst school began to fall out of favor among Homeric scholars.[20] It did not die out entirely, but it came to be increasingly seen as a discredited dead end.[20] Starting in around 1928, Milman Parry and Albert Lord, after their studies of folk bards in the Balkans, developed the "Oral-Formulaic Theory" that the Homeric poems were originally composed through improvised oral performances, which relied on traditional epithets and poetic formulas.[28][27][20] This theory found very wide scholarly acceptance[28][27][20] and explained many previously puzzling features of the Homeric poems, including their unusually archaic language, their extensive use of stock epithets, and their other "repetitive" features.[27] Many scholars concluded that the "Homeric question" had finally been answered.[20] Meanwhile, the 'Neoanalysts' sought to bridge the gap between the 'Analysts' and 'Unitarians'.[29][30] The Neoanalysts sought to trace the relationships between the Homeric poems and other epic poems, which have now been lost, but of which modern scholars do possess some patchy knowledge.[20] Knowledge of earlier versions of the epics can be derived from anomalies of structure and detail in our surviving version of the Iliad and Odyssey.  These anomalies point to earlier versions of the Iliad in which Ajax played a more prominent role, in which the Achaean embassy to Achilles comprised different characters, and in which Patroclus was actually mistaken for Achilles by the Trojans.  They point to earlier versions of the Odyssey in which Telemachus went in search of news of his father not to Menelaus in Sparta but to Idomeneus in Crete, in which Telemachus met up with his father in Crete and conspired with him to return to Ithaca disguised as the soothsayer Theoclymenus, and in which Penelope recognized Odysseus much earlier in the narrative and conspired with him in the destruction of the suitors.[31]  Neoanalysis can be viewed as a form of Analysis informed by the principles of Oral Theory, recognizing as it does the existence and influence of previously existing tales and yet appreciating the technique of a single poet in adapting them to his Iliad and Odyssey.
+
+Most contemporary scholars, although they disagree on other questions about the genesis of the poems, agree that the Iliad and the Odyssey were not produced by the same author, based on "the many differences of narrative manner, theology, ethics, vocabulary, and geographical perspective, and by the apparently imitative character of certain passages of the Odyssey in relation to the Iliad."[32][33][34][20] Nearly all scholars agree that the Iliad and the Odyssey are unified poems, in that each poem shows a clear overall design, and that they are not merely strung together from unrelated songs.[20] It is also generally agreed that each poem was composed mostly by a single author, who probably relied heavily on older oral traditions.[20] Nearly all scholars agree that the Doloneia in Book X of the Iliad is not part of the original poem, but rather a later insertion by a different poet.[20]
+
+Some ancient scholars believed Homer to have been an eyewitness to the Trojan War; others thought he had lived up to 500 years afterwards.[35] Contemporary scholars continue to debate the date of the poems.[36][37][20] A long history of oral transmission lies behind the composition of the poems, complicating the search for a precise date.[38] At one extreme, Richard Janko has proposed a date for both poems to the eighth century BC based on linguistic analysis and statistics.[36][37] Barry B. Powell dates the composition of the Iliad and the Odyssey to sometime between 800 and 750 BC, based on the statement from Herodotus, who lived in the late fifth century BC, that Homer lived four hundred years before his own time "and not more" (καὶ οὐ πλέοσι), and on the fact that the poems do not mention hoplite battle tactics, inhumation, or literacy.[39] Martin Litchfield West has argued that the Iliad echoes the poetry of Hesiod, and that it must have been composed around 660–650 BC at the earliest, with the Odyssey up to a generation later.[40][41][20] He also interprets passages in the Iliad as showing knowledge of historical events that occurred in the ancient Near East during the middle of the seventh century BC, including the destruction of Babylon by Sennacherib in 689 BC and the Sack of Thebes by Ashurbanipal in 663/4 BC.[20] At the other extreme, a few American scholars such as Gregory Nagy see "Homer" as a continually evolving tradition, which grew much more stable as the tradition progressed, but which did not fully cease to continue changing and evolving until as late as the middle of the second century BC.[36][37][20]
+
+"'Homer" is a name of unknown etymological origin, around which many theories were erected in antiquity. One such linkage was to the Greek ὅμηρος (hómēros), "hostage" (or "surety"). The explanations suggested by modern scholars tend to mirror their position on the overall Homeric question. Nagy interprets it as "he who fits (the song) together". West has advanced both possible Greek and Phoenician etymologies.[42][43]
+
+Scholars continue to debate questions such as whether the Trojan War actually took place – and if so when and where – and to what extent the society depicted by Homer is based on his own or one which was, even at the time of the poems' composition, known only as legend. The Homeric epics are largely set in the east and center of the Mediterranean, with some scattered references to Egypt, Ethiopia and other distant lands, in a warlike society that resembles that of the Greek world slightly before the hypothesized date of the poems' composition.[44][45][46][47]
+
+In ancient Greek chronology, the sack of Troy was dated to 1184 BC. By the nineteenth century, there was widespread scholarly skepticism that the Trojan War had ever happened and that Troy had even existed, but in 1873 Heinrich Schliemann announced to the world that he had discovered the ruins of Homer's Troy at Hissarlik in modern Turkey. Some contemporary scholars think the destruction of Troy VIIa circa 1220 BC was the origin of the myth of the Trojan War, others that the poem was inspired by multiple similar sieges that took place over the centuries.[48]
+
+Most scholars now agree that the Homeric poems depict customs and elements of the material world that are derived from different periods of Greek history.[27][49][50] For instance, the heroes in the poems use bronze weapons, characteristic of the Bronze Age in which the poems are set, rather than the later Iron Age during which they were composed;[27][49][50] yet the same heroes are cremated (an Iron Age practice) rather than buried (as they were in the Bronze Age).[27][49][50] In some parts of the Homeric poems, heroes are accurately described as carrying large shields like those used by warriors during the Mycenaean period,[27] but, in other places, they are instead described carrying the smaller shields that were commonly used during the time when the poems were written in the early Iron Age.[27]
+
+In the Iliad 10.260–265, Odysseus is described as wearing a helmet made of boar's tusks. Such helmets were not worn in Homer's time, but were commonly worn by aristocratic warriors between 1600 and 1150 BC.[51][52][53] The decipherment of Linear B in the 1950s by Michael Ventris and continued archaeological investigation has increased modern scholars' understanding of Aegean civilisation, which in many ways resembles the ancient Near East more than the society described by Homer.[54] Some aspects of the Homeric world are simply made up;[27] for instance, the Iliad 22.145–56 describes there being two springs that run near the city of Troy, one that runs steaming hot and the other that runs icy cold.[27] It is here that Hector takes his final stand against Achilles.[27] Archaeologists, however, have uncovered no evidence that springs of this description ever actually existed.[27]
+
+The Homeric epics are written in an artificial literary language or 'Kunstsprache' only used in epic hexameter poetry. Homeric Greek shows features of multiple regional Greek dialects and periods, but is fundamentally based on Ionic Greek, in keeping with the tradition that Homer was from Ionia. Linguistic analysis suggests that the Iliad was composed slightly before the Odyssey, and that Homeric formulae preserve older features than other parts of the poems.[55][56]
+
+The Homeric poems were composed in unrhymed dactylic hexameter; ancient Greek metre was quantity-based rather than stress-based.[57][58] Homer frequently uses set phrases such as epithets ('crafty Odysseus', 'rosy-fingered Dawn', 'owl-eyed Athena', etc.), Homeric formulae ('and then answered [him/her], Agamemnon, king of men', 'when the early-born rose-fingered Dawn came to light', 'thus he/she spoke'), simile, type scenes, ring composition and repetition. These habits aid the extemporizing bard, and are characteristic of oral poetry. For instance, the main words of a Homeric sentence are generally placed towards the beginning, whereas literate poets like Virgil or Milton use longer and more complicated syntactical structures. Homer then expands on these ideas in subsequent clauses; this technique is called parataxis.[59]
+
+The so-called 'type scenes' (typischen Scenen), were named by Walter Arend in 1933. He noted that Homer often, when describing frequently recurring activities such as eating, praying, fighting and dressing, used blocks of set phrases in sequence that were then elaborated by the poet. The 'Analyst' school had considered these repetitions as un-Homeric, whereas Arend interpreted them philosophically. Parry and Lord noted that these conventions are found in many other cultures.[60][61]
+
+'Ring composition' or chiastic structure (when a phrase or idea is repeated at both the beginning and end of a story, or a series of such ideas first appears in the order A, B, C... before being reversed as ...C, B, A) has been observed in the Homeric epics. Opinion differs as to whether these occurrences are a conscious artistic device, a mnemonic aid or a spontaneous feature of human storytelling.[62][63]
+
+Both of the Homeric poems begin with an invocation to the Muse.[64] In the Iliad, the poet invokes her to sing of "the anger of Achilles",[64] and, in the Odyssey, he asks her to sing of "the man of many ways".[64] A similar opening was later employed by Virgil in his Aeneid.[64]
+
+The orally transmitted Homeric poems were put into written form at some point between the eighth and sixth centuries BC. Some scholars believe that they were dictated to a scribe by the poet and that our inherited versions of the Iliad and Odyssey were in origin orally-dictated texts.[65] Albert Lord noted that the Balkan bards that he was studying revised and expanded their songs in their process of dictating.[66] Some scholars hypothesize that a similar process of revision and expansion occurred when the Homeric poems were first written down.[67][68]  Other scholars hold that, after the poems were created in the eighth century, they continued to be orally transmitted with considerable revision until they were written down in the sixth century.[69] After textualisation, the poems were each divided into 24 rhapsodes, today referred to as books, and labelled by the letters of the Greek alphabet. Most scholars attribute the book divisions to the Hellenistic scholars of Alexandria, in Egypt.[70]  Some trace the divisions back further to the Classical period.[71]  Very few credit Homer himself with the divisions.[72]
+
+In antiquity, it was widely held that the Homeric poems were collected and organised in Athens in the late sixth century BC by the tyrant Peisistratos (died 528/7 BC), in what subsequent scholars have dubbed the "Peisistratean recension".[73][21] The idea that the Homeric poems were originally transmitted orally and first written down during the reign of Peisistratos is referenced by the first-century BC Roman orator Cicero and is also referenced in a number of other surviving sources, including two ancient Lives of Homer.[21] From around 150 BC, the texts of the Homeric poems seem to have become relatively established. After the establishment of the Library of Alexandria, Homeric scholars such as Zenodotus of Ephesus, Aristophanes of Byzantium and in particular Aristarchus of Samothrace helped establish a canonical text.[74]
+
+The first printed edition of Homer was produced in 1488 in Milan, Italy. Today scholars use medieval manuscripts, papyri and other sources; some argue for a "multi-text" view, rather than seeking a single definitive text. The nineteenth-century edition of Arthur Ludwich mainly follows Aristarchus's work, whereas van Thiel's (1991, 1996) follows the medieval vulgate. Others, such as Martin West (1998–2000) or T.W. Allen, fall somewhere between these two extremes.[74]
+
+This is a partial list of translations into English of Homer's Iliad and Odyssey.

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+The equator of a rotating spheroid (such as a planet) is the parallel (circle of latitude) at which latitude is defined to be 0°. It is the imaginary line on the spheroid, equidistant from its poles, dividing it into northern and southern hemispheres. In other words, it is the intersection of the spheroid with the plane perpendicular to its axis of rotation and midway between its geographical poles.
+
+On Earth, the equator is about 40,075 km (24,901 mi) long, of which 78.8% lies across water and 21.3% over land. Indonesia is the country straddling the greatest length of the equatorial line across both land and sea.
+
+The name is derived from medieval Latin word aequator, in the phrase circulus aequator diei et noctis, meaning 'circle equalizing day and night', from the Latin word aequare meaning 'make equal'.[1]
+
+The latitude of the Earth's equator is, by definition, 0° (zero degrees) of arc. The equator is one of the five notable circles of latitude on Earth; the other four are both polar circles (the Arctic Circle and the Antarctic Circle) and both tropical circles (the Tropic of Cancer and the Tropic of Capricorn). The equator is the only line of latitude which is also a great circle—that is, one whose plane passes through the center of the globe. The plane of Earth's equator, when projected outwards to the celestial sphere, defines the celestial equator.
+
+In the cycle of Earth's seasons, the equatorial plane runs through the Sun twice per year: on the equinoxes in March and September. To a person on Earth, the Sun appears to travel above the equator (or along the celestial equator) at these times. Light rays from the Sun's center are perpendicular to Earth's surface at the point of solar noon on the equator.
+
+Locations on the equator experience the shortest sunrises and sunsets because the Sun's daily path is nearly perpendicular to the horizon for most of the year. The length of daylight (sunrise to sunset) is almost constant throughout the year; it is about 14 minutes longer than nighttime due to atmospheric refraction and the fact that sunrise begins (or sunset ends) as the upper limb, not the center, of the Sun's disk contacts the horizon.
+
+Earth bulges slightly at the equator; the "average" diameter of Earth is 12,750 km (7,920 mi), but the diameter at the equator is about 43 km (27 mi) greater than at the poles.[2]
+
+Sites near the equator, such as the Guiana Space Centre in Kourou, French Guiana, are good locations for spaceports as they have a fastest rotational speed of any latitude, 460 m/s. The added velocity reduces the fuel needed to launch spacecraft eastward (in the direction of Earth's rotation) to orbit, while simultaneously avoiding costly maneuvers to flatten inclination during missions such as the Apollo moon landings.[3]
+
+The precise location of the equator is not truly fixed; the true equatorial plane is perpendicular to the Earth's spin axis, which drifts about 9 metres (30 ft) during a year. This effect must be accounted for in detailed geophysical measurements.[citation needed]
+
+Geological samples show the equator significantly changed positions between 12 to 48 million years ago, as sediment deposited by ocean thermal currents at the equator have shifted. The deposits by thermal currents are determined by the axis of the earth, which determines solar coverage of the Earth’s surface. Changes in Earth axis can also be observed in the geographic layout of volcanic island chains, which are created by shifting hot spots under the Earth’s crust as the axis and crust move.[4]
+
+The International Association of Geodesy (IAG) and the International Astronomical Union (IAU) have chosen to use an equatorial radius of 6,378.1366 kilometres (3,963.1903 mi) (codified as the IAU 2009 value).[5] This equatorial radius is also in the 2003 and 2010 IERS Conventions.[6] It is also the equatorial radius used for the IERS 2003 ellipsoid.  If it were really circular, the length of the equator would then be exactly 2π times the radius, namely 40,075.0142 kilometres (24,901.4594 mi). The GRS 80 (Geodetic Reference System 1980) as approved and adopted by the IUGG at its Canberra, Australia meeting of 1979 has an equatorial radius of 6,378.137 kilometres (3,963.191 mi). The WGS 84 (World Geodetic System 1984) which is a standard for use in cartography, geodesy, and satellite navigation including GPS, also has an equatorial radius of 6,378.137 kilometres (3,963.191 mi). For both GRS 80 and WGS 84, this results in a length for the equator of 40,075.0167 km (24,901.4609 mi).
+
+The geographical mile is defined as one arc-minute of the equator, so it has different values depending on which radius is assumed. For example, by WSG-84, the distance is 1,855.3248 metres (6,087.024 ft), while by IAU-2000, it is 1,855.3257 metres (6,087.027 ft). This is a difference of less than one millimetre (0.039 in) over the total distance (approximately 1.86 kilometres or 1.16 miles).
+
+The earth is commonly modeled as a sphere flattened 0.336% along its axis. This makes the equator 0.16% longer than a meridian (a great circle passing through the two poles). The IUGG standard meridian is, to the nearest millimetre, 40,007.862917 kilometres (24,859.733480 mi), one arc-minute of which is 1,852.216 metres (6,076.82 ft), explaining the SI standardization of the nautical mile as 1,852 metres (6,076 ft), more than 3 metres (9.8 ft) less than the geographical mile.
+
+The sea-level surface of the Earth (the geoid) is irregular, so the actual length of the equator is not so easy to determine. Aviation Week and Space Technology on 9 October 1961 reported that measurements using the Transit IV-A satellite had shown the equatorial diameter from longitude 11° West to 169° East to be 1,000 feet (300 m) greater than its diameter ninety degrees away.[citation needed]
+
+The equator passes through the land of 11 countries. Starting at the Prime Meridian and heading eastwards, the equator passes through:
+
+Despite its name, no part of Equatorial Guinea lies on the equator. However, its island of Annobón is 155 km (96 mi) south of the equator, and the rest of the country lies to the north.
+
+Seasons result from the tilt of the Earth's axis compared to the plane of its revolution around the Sun. Throughout the year the northern and southern hemispheres are alternately turned either toward or away from the sun depending on Earth's position in its orbit. The hemisphere turned toward the sun receives more sunlight and is in summer, while the other hemisphere receives less sun and is in winter (see solstice).
+
+At the equinoxes, the Earth's axis is perpendicular to the sun rather than tilted toward or away, meaning that day and night are both about 12 hours long across the whole of the Earth.
+
+Near the equator, this means the variation in strength of solar radiation is different relative to the time of year than it is at higher latitudes: Maximum solar radiation is received during the equinoxes, when a place at the equator is under the subsolar point at high noon, and the intermediate seasons of spring and autumn occur at higher latitudes, and the minimum occurs during both solstices, when either pole is tilted towards or away from the sun, resulting in either summer or winter in both hemispheres. This also results in a corresponding movement of the equator away from the subsolar point, which is then situated over or near the relevant tropic circle. Nevertheless, temperatures are high year round due to the earth's axial tilt of 23.5° not being enough to create a low minimum midday declination to sufficiently weaken the sun's rays even during the solstices.
+
+Near the equator there is little temperature change throughout the year, though there may be dramatic differences in rainfall and humidity. The terms summer, autumn, winter and spring do not generally apply. Lowlands around the equator generally have a tropical rainforest climate, also known as an equatorial climate, though cold ocean currents cause some regions to have tropical monsoon climates with a dry season in the middle of the year, and the Somali Current generated by the Asian monsoon due to continental heating via the high Tibetan Plateau causes Greater Somalia to have an arid climate despite its equatorial location.
+
+Average annual temperatures in equatorial lowlands are around 31 °C (88 °F) during the afternoon and 23 °C (73 °F) around sunrise. Rainfall is very high away from cold ocean current upwelling zones, from 2,500 to 3,500 mm (100 to 140 in) per year. There are about 200 rainy days per year and average annual sunshine hours are around 2,000. Despite high year-round sea level temperatures, some higher altitudes such as the Andes and Mount Kilimanjaro have glaciers. The highest point on the equator is at the elevation of 4,690 metres (15,387 ft), at 0°0′0″N 77°59′31″W / 0.00000°N 77.99194°W / 0.00000; -77.99194 (highest point on the equator), found on the southern slopes of Volcán Cayambe [summit 5,790 metres (18,996 ft)] in Ecuador. This is slightly above the snow line and is the only place on the equator where snow lies on the ground. At the equator, the snow line is around 1,000 metres (3,300 ft) lower than on Mount Everest and as much as 2,000 metres (6,600 ft) lower than the highest snow line in the world, near the Tropic of Capricorn on Llullaillaco.
+
+There is a widespread maritime tradition of holding ceremonies to mark a sailor's first crossing of the equator. In the past, these ceremonies have been notorious for their brutality, especially in naval practice.[citation needed] Milder line-crossing ceremonies, typically featuring King Neptune, are also held for passengers' entertainment on some civilian ocean liners and cruise ships.[citation needed]

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+Coordinates: 2°00′S 77°30′W / 2.000°S 77.500°W / -2.000; -77.500
+
+in South America (grey)
+
+Ecuador (/ˈɛkwədɔːr/ (listen) EK-wə-dor; Spanish pronunciation: [ekwaˈðoɾ] (listen); Quechua: Ikwayur; Shuar: Ecuador or Ekuatur.),[12][13] officially the Republic of Ecuador (Spanish: República del Ecuador,  which literally translates as "Republic of the Equator"; Quechua: Ikwadur Ripuwlika; Shuar: Ekuatur Nunka),[14][15] is a country in northwestern South America, bordered by Colombia on the north, Peru on the east and south, and the Pacific Ocean on the west. Ecuador also includes the Galápagos Islands in the Pacific, about 1,000 kilometres (621 mi) west of the mainland. The capital is Quito.[16][17]
+
+The territories of modern-day Ecuador were once home to a variety of Amerindian groups that were gradually incorporated into the Inca Empire during the 15th century. The territory was colonized by Spain during the 16th century, achieving independence in 1820 as part of Gran Colombia, from which it emerged as its own sovereign state in 1830. The legacy of both empires is reflected in Ecuador's ethnically diverse population, with most of its 17.1 million people being mestizos, followed by large minorities of European, Amerindian, and African descendants. Spanish is the official language and is spoken by a majority of the population, though 13 Amerindian languages are also recognized, including Quechua and Shuar.
+
+The sovereign state of Ecuador is a middle-income representative democratic republic and a developing country[18] that is highly dependent on commodities, namely petroleum and agricultural products. It is governed as a democratic presidential republic. One of 17 megadiverse countries in the world,[19][20] Ecuador hosts many endemic plants and animals, such as those of the Galápagos Islands. In recognition of its unique ecological heritage, the new constitution of 2008 is the first in the world to recognize legally enforceable Rights of Nature, or ecosystem rights.[21] It also has the fifth lowest homicide rate in the Americas.[22] Between 2006 and 2016, poverty decreased from 36.7% to 22.5% and annual per capita GDP growth was 1.5 percent (as compared to 0.6 percent over the prior two decades). At the same time, the country's Gini index of economic inequality decreased from 0.55 to 0.47.[23]
+
+Various peoples had settled in the area of future Ecuador before the arrival of the Incas. The archeological evidence suggests that the Paleo-Indians' first dispersal into the Americas occurred near the end of the last glacial period, around 16,500–13,000 years ago. The first Indians who reached Ecuador may have journeyed by land from North and Central America or by boat down the Pacific Ocean coastline. Much later migrations to Ecuador may have come via the Amazon tributaries, others descended from northern South America, and others ascended from the southern part of South America through the Andes. They developed different languages while emerging as unique ethnic groups.
+
+Even though their languages were unrelated, these groups developed similar groups of cultures, each based in different environments. The people of the coast developed a fishing, hunting, and gathering culture; the people of the highland Andes developed a sedentary agricultural way of life, and the people of the Amazon basin developed a nomadic hunting-and-gathering mode of existence.
+
+Over time these groups began to interact and intermingle with each other so that groups of families in one area became one community or tribe, with a similar language and culture. Many civilizations arose in Ecuador, such as the Valdivia Culture and Machalilla Culture on the coast, the Quitus (near present-day Quito), and the Cañari (near present-day Cuenca). Each civilisation developed its own distinctive architecture, pottery, and religious interests.
+
+In the highland Andes mountains, where life was more sedentary, groups of tribes cooperated and formed villages; thus the first nations based on agricultural resources and the domestication of animals formed. Eventually, through wars and marriage alliances of their leaders, a group of nations formed confederations. One region consolidated under a confederation called the Shyris, which exercised organized trading and bartering between the different regions. Its political and military power came under the rule of the Duchicela blood-line.
+
+When the Incas arrived, they found that these confederations were so developed that it took the Incas two generations of rulers—Topa Inca Yupanqui and Huayna Capac—to absorb them into the Inca Empire. The native confederations that gave them the most problems were deported to distant areas of Peru, Bolivia, and north Argentina. Similarly, a number of loyal Inca subjects from Peru and Bolivia were brought to Ecuador to prevent rebellion. Thus, the region of highland Ecuador became part of the Inca Empire in 1463 sharing the same language.
+
+In contrast, when the Incas made incursions into coastal Ecuador and the eastern Amazon jungles of Ecuador, they found both the environment and indigenous people more hostile. Moreover, when the Incas tried to subdue them, these indigenous people withdrew to the interior and resorted to guerrilla tactics. As a result, Inca expansion into the Amazon Basin and the Pacific coast of Ecuador was hampered. The indigenous people of the Amazon jungle and coastal Ecuador remained relatively autonomous until the Spanish soldiers and missionaries arrived in force. The Amazonian people and the Cayapas of Coastal Ecuador were the only groups to resist Inca and Spanish domination, maintaining their language and culture well into the 21st century.
+
+Before the arrival of the Spaniards, the Inca Empire was involved in a civil war. The untimely death of both the heir Ninan Cuchi and the Emperor Huayna Capac, from a European disease that spread into Ecuador, created a power vacuum between two factions. The northern faction headed by Atahualpa claimed that Huayna Capac gave a verbal decree before his death about how the empire should be divided. He gave the territories pertaining to present-day Ecuador and northern Peru to his favorite son Atahualpa, who was to rule from Quito; and he gave the rest to Huáscar, who was to rule from Cuzco. He willed that his heart be buried in Quito, his favorite city, and the rest of his body be buried with his ancestors in Cuzco.
+
+Huáscar did not recognize his father's will, since it did not follow Inca traditions of naming an Inca through the priests. Huáscar ordered Atahualpa to attend their father's burial in Cuzco and pay homage to him as the new Inca ruler. Atahualpa, with a large number of his father's veteran soldiers, decided to ignore Huáscar, and a civil war ensued. A number of bloody battles took place until finally Huáscar was captured. Atahualpa marched south to Cuzco and massacred the royal family associated with his brother.
+
+In 1532, a small band of Spaniards headed by Francisco Pizarro landed in Tumbez and marched over the Andes Mountains until they reached Cajamarca, where the new Inca Atahualpa was to hold an interview with them. Valverde, the priest, tried to convince Atahualpa that he should join the Catholic Church and declare himself a vassal of Spain. This infuriated Atahualpa so much that he threw the Bible to the ground. At this point the enraged Spaniards, with orders from Valverde, attacked and massacred unarmed escorts of the Inca and captured Atahualpa. Pizarro promised to release Atahualpa if he made good his promise of filling a room full of gold. But, after a mock trial, the Spaniards executed Atahualpa by strangulation.
+
+New infectious diseases such as smallpox, endemic to the Europeans, caused high fatalities among the Amerindian population during the first decades of Spanish rule, as they had no immunity. At the same time, the natives were forced into the encomienda labor system for the Spanish. In 1563, Quito became the seat of a real audiencia (administrative district) of Spain and part of the Viceroyalty of Peru and later the Viceroyalty of New Granada.
+
+The 1797 Riobamba earthquake, which caused up to 40,000 casualties, was studied by Alexander von Humboldt, when he visited the area in 1801–1802.[24]
+
+After nearly 300 years of Spanish rule, Quito was still a small city numbering 10,000 inhabitants. On August 10, 1809, the city's criollos called for independence from Spain (first among the peoples of Latin America). They were led by Juan Pío Montúfar, Quiroga, Salinas, and Bishop Cuero y Caicedo. Quito's nickname, "Luz de América" ("Light of America"), is based on its leading role in trying to secure an independent, local government. Although the new government lasted no more than two months, it had important repercussions and was an inspiration for the independence movement of the rest of Spanish America. August 10 is now celebrated as Independence Day, a national holiday.[25]
+
+On October 9, 1820, the Department of Guayaquil became the first territory in Ecuador to gain its independence from Spain, and it spawned most of the Ecuadorian coastal provinces, establishing itself as an independent state. Its inhabitants celebrated what is now Ecuador's official Independence Day on May 24, 1822. The rest of Ecuador gained its independence after Antonio José de Sucre defeated the Spanish Royalist forces at the Battle of Pichincha, near Quito. Following the battle, Ecuador joined Simón Bolívar's Republic of Gran Colombia, also including modern-day Colombia, Venezuela and Panama. In 1830, Ecuador separated from Gran Colombia and became an independent republic.
+
+The 19th century was marked by instability for Ecuador with a rapid succession of rulers. The first president of Ecuador was the Venezuelan-born Juan José Flores, who was ultimately deposed, followed by several authoritarian leaders, such as Vicente Rocafuerte; José Joaquín de Olmedo; José María Urbina; Diego Noboa; Pedro José de Arteta; Manuel de Ascásubi; and Flores's own son, Antonio Flores Jijón, among others. The conservative Gabriel Garcia Moreno unified the country in the 1860s with the support of the Roman Catholic Church. In the late 19th century, world demand for cocoa tied the economy to commodity exports and led to migrations from the highlands to the agricultural frontier on the coast.
+
+Ecuador abolished slavery and freed its black slaves in 1851.[26]
+
+The Liberal Revolution of 1895 under Eloy Alfaro reduced the power of the clergy and the conservative land owners. This liberal wing retained power until the military "Julian Revolution" of 1925. The 1930s and 1940s were marked by instability and emergence of populist politicians, such as five-time President José María Velasco Ibarra.
+
+Brasilia Presidential Act
+
+Since Ecuador's separation from Colombia on May 13, 1830, its first President, General Juan José Flores, laid claim to the territory that was called the Real Audiencia of Quito, also referred to as the Presidencia of Quito. He supported his claims with Spanish Royal decrees or Real Cedulas, that delineated the borders of Spain's former overseas colonies. In the case of Ecuador, Flores-based Ecuador's de jure claims on the following cedulas - Real Cedula of 1563, 1739, and 1740; with modifications in the Amazon Basin and Andes Mountains that were introduced through the Treaty of Guayaquil (1829) which Peru reluctantly signed, after the overwhelmingly outnumbered Gran Colombian force led by Antonio José de Sucre defeated President and General La Mar's Peruvian invasion force in the Battle of Tarqui. In addition, Ecuador's eastern border with the Portuguese colony of Brazil in the Amazon Basin was modified before the wars of Independence by the First Treaty of San Ildefonso (1777) between the Spanish Empire and the Portuguese Empire. Moreover, to add legitimacy to his claims, on February 16, 1840, Flores signed a treaty with Spain, whereby Flores convinced Spain to officially recognize Ecuadorian independence and its sole rights to colonial titles over Spain's former colonial territory known anciently to Spain as the Kingdom and Presidency of Quito.
+
+Ecuador during its long and turbulent history has lost most of its contested territories to each of its more powerful neighbors, such as Colombia in 1832 and 1916, Brazil in 1904 through a series of peaceful treaties, and Peru after a short war in which the Protocol of Rio de Janeiro was signed in 1942.
+
+During the struggle for independence, before Peru or Ecuador became independent nations, a few areas of the former Vice Royalty of New Granada - Guayaquil, Tumbez, and Jaén - declared themselves independent from Spain. A few months later, a part of the Peruvian liberation army of San Martin decided to occupy the independent cities of Tumbez and Jaén with the intention of using these towns as springboards to occupy the independent city of Guayaquil and then to liberate the rest of the Audiencia de Quito (Ecuador). It was common knowledge among the top officers of the liberation army from the south that their leader San Martin wished to liberate present-day Ecuador and add it to the future republic of Peru, since it had been part of the Inca Empire before the Spaniards conquered it.
+
+However, Bolívar's intention was to form a new republic known as the Gran Colombia, out of the liberated Spanish territory of New Granada which consisted of Colombia, Venezuela, and Ecuador. San Martin's plans were thwarted when Bolívar, with the help of Marshal Antonio José de Sucre and the Gran Colombian liberation force, descended from the Andes mountains and occupied Guayaquil; they also annexed the newly liberated Audiencia de Quito to the Republic of Gran Colombia. This happened a few days before San Martin's Peruvian forces could arrive and occupy Guayaquil, with the intention of annexing Guayaquil to the rest of Audiencia of Quito (Ecuador) and to the future republic of Peru. Historic documents repeatedly stated that San Martin told Bolivar he came to Guayaquil to liberate the land of the Incas from Spain. Bolivar countered by sending a message from Guayaquil welcoming San Martin and his troops to Colombian soil.
+
+In the south, Ecuador had de jure claims to a small piece of land beside the Pacific Ocean known as Tumbes which lay between the Zarumilla and Tumbes rivers. In Ecuador's southern Andes Mountain region where the Marañon cuts across, Ecuador had de jure claims to an area it called Jaén de Bracamoros. These areas were included as part of the territory of Gran Colombia by Bolivar on December 17, 1819, during the Congress of Angostura when the Republic of Gran Colombia was created. Tumbes declared itself independent from Spain on January 17, 1821, and Jaen de Bracamoros on June 17, 1821, without any outside help from revolutionary armies. However, that same year, 1821, Peruvian forces participating in the Trujillo revolution occupied both Jaen and Tumbes. Some Peruvian generals, without any legal titles backing them up and with Ecuador still federated with the Gran Colombia, had the desire to annex Ecuador to the Republic of Peru at the expense of the Gran Colombia, feeling that Ecuador was once part of the Inca Empire.
+
+On July 28, 1821, Peruvian independence was proclaimed in Lima by the Liberator San Martin, and Tumbes and Jaen, which were included as part of the revolution of Trujillo by the Peruvian occupying force, had the whole region swear allegiance to the new Peruvian flag and incorporated itself into Peru, even though Peru was not completely liberated from Spain. After Peru was completely liberated from Spain by the patriot armies led by Bolivar and Antonio Jose de Sucre at the Battle of Ayacucho dated December 9, 1824, there was a strong desire by some Peruvians to resurrect the Inca Empire and to include Bolivia and Ecuador. One of these Peruvian Generals was the Ecuadorian-born José de La Mar, who became one of Peru's presidents after Bolivar resigned as dictator of Peru and returned to Colombia. Gran Colombia had always protested Peru for the return of Jaen and Tumbes for almost a decade, then finally Bolivar after long and futile discussion over the return of Jaen, Tumbes, and part of Mainas, declared war. President and General José de La Mar, who was born in Ecuador, believing his opportunity had come to annex the District of Ecuador to Peru, personally, with a Peruvian force, invaded and occupied Guayaquil and a few cities in the Loja region of southern Ecuador on November 28, 1828.
+
+The war ended when a triumphant heavily outnumbered southern Gran Colombian army at Battle of Tarqui dated February 27, 1829, led by Antonio José de Sucre, defeated the Peruvian invasion force led by President La Mar. This defeat led to the signing of the Treaty of Guayaquil dated September 22, 1829, whereby Peru and its Congress recognized Gran Colombian rights over Tumbes, Jaen, and Maynas. Through protocolized meetings between representatives of Peru and Gran Colombia, the border was set as Tumbes river in the west and in the east the Maranon and Amazon rivers were to be followed toward Brazil as the most natural borders between them. However, what was pending was whether the new border around the Jaen region should follow the Chinchipe River or the Huancabamba River. According to the peace negotiations Peru agreed to return Guayaquil, Tumbez, and Jaén; despite this, Peru returned Guayaquil, but failed to return Tumbes and Jaén, alleging that it was not obligated to follow the agreements, since the Gran Colombia ceased to exist when it divided itself into three different nations - Ecuador, Colombia, and Venezuela.
+
+The Central District of the Gran Colombia, known as Cundinamarca or New Granada (modern Colombia) with its capital in Bogota, did not recognize the separation of the Southern District of the Gran Colombia, with its capital in Quito, from the Gran Colombian federation on May 13, 1830. After Ecuador's separation, the Department of Cauca voluntarily decided to unite itself with Ecuador due to instability in the central government of Bogota. The Venezuelan born President of Ecuador, the general Juan José Flores, with the approval of the Ecuadorian congress annexed the Department of Cauca on December 20, 1830, since the government of Cauca had called for union with the District of the South as far back as April 1830. Moreover, the Cauca region, throughout its long history, had very strong economic and cultural ties with the people of Ecuador. Also, the Cauca region, which included such cities as Pasto, Popayán, and Buenaventura, had always been dependent on the Presidencia or Audiencia of Quito.
+
+Fruitless negotiations continued between the governments of Bogotá and Quito, where the government of Bogotá did not recognize the separation of Ecuador or that of Cauca from the Gran Colombia until war broke out in May 1832. In five months, New Granada defeated Ecuador due to the fact that the majority of the Ecuadorian Armed Forces were composed of rebellious angry unpaid veterans from Venezuela and Colombia that did not want to fight against their fellow countrymen. Seeing that his officers were rebelling, mutinying, and changing sides, President Flores had no option but to reluctantly make peace with New Granada. The Treaty of Pasto of 1832 was signed by which the Department of Cauca was turned over to New Granada (modern Colombia), the government of Bogotá recognized Ecuador as an independent country and the border was to follow the Ley de División Territorial de la República de Colombia (Law of the Division of Territory of the Gran Colombia) passed on June 25, 1824. This law set the border at the river Carchi and the eastern border that stretched to Brazil at the Caquetá river. Later, Ecuador contended that the Republic of Colombia, while reorganizing its government, unlawfully made its eastern border provisional and that Colombia extended its claims south to the Napo River because it said that the Government of Popayán extended its control all the way to the Napo River.
+
+When Ecuador seceded from the Gran Colombia, Peru decided not to follow the treaty of Guayaquil of 1829 or the protocoled agreements made. Peru contested Ecuador's claims with the newly discovered Real Cedula of 1802, by which Peru claims the King of Spain had transferred these lands from the Viceroyalty of New Granada to the Viceroyalty of Peru. During colonial times this was to halt the ever-expanding Portuguese settlements into Spanish domains, which were left vacant and in disorder after the expulsion of Jesuit missionaries from their bases along the Amazon Basin. Ecuador countered by labeling the Cedula of 1802 an ecclesiastical instrument, which had nothing to do with political borders. Peru began its de facto occupation of disputed Amazonian territories, after it signed a secret 1851 peace treaty in favor of Brazil. This treaty disregarded Spanish rights that were confirmed during colonial times by a Spanish-Portuguese treaty over the Amazon regarding territories held by illegal Portuguese settlers.
+
+Peru began occupying the defenseless missionary villages in the Mainas or Maynas region, which it began calling Loreto, with its capital in Iquitos. During its negotiations with Brazil, Peru stated that based on the royal cedula of 1802, it claimed Amazonian Basin territories up to Caqueta River in the north and toward the Andes Mountain range, depriving Ecuador and Colombia of all their claims to the Amazon Basin. Colombia protested stating that its claims extended south toward the Napo and Amazon Rivers. Ecuador protested that it claimed the Amazon Basin between the Caqueta river and the Marañon-Amazon river. Peru ignored these protests and created the Department of Loreto in 1853 with its capital in Iquitos which it had recently invaded and systematically began to occupy using the river systems in all the territories claimed by both Colombia and Ecuador. Peru briefly occupied Guayaquil again in 1860, since Peru thought that Ecuador was selling some of the disputed land for development to British bond holders, but returned Guayaquil after a few months. The border dispute was then submitted to Spain for arbitration from 1880 to 1910, but to no avail.
+
+In the early part of the 20th century, Ecuador made an effort to peacefully define its eastern Amazonian borders with its neighbours through negotiation. On May 6, 1904, Ecuador signed the Tobar-Rio Branco Treaty recognizing Brazil's claims to the Amazon in recognition of Ecuador's claim to be an Amazonian country to counter Peru's earlier Treaty with Brazil back on October 23, 1851. Then after a few meetings with the Colombian government's representatives an agreement was reached and the Muñoz Vernaza-Suarez Treaty was signed July 15, 1916, in which Colombian rights to the Putumayo river were recognized as well as Ecuador's rights to the Napo river and the new border was a line that ran midpoint between those two rivers. In this way, Ecuador gave up the claims it had to the Amazonian territories between the Caquetá River and Napo River to Colombia, thus cutting itself off from Brazil. Later, a brief war erupted between Colombia and Peru, over Peru's claims to the Caquetá region, which ended with Peru reluctantly signing the Salomon-Lozano Treaty on March 24, 1922. Ecuador protested this secret treaty, since Colombia gave away Ecuadorian claimed land to Peru that Ecuador had given to Colombia in 1916.
+
+On July 21, 1924, the Ponce-Castro Oyanguren Protocol was signed between Ecuador and Peru where both agreed to hold direct negotiations and to resolve the dispute in an equitable manner and to submit the differing points of the dispute to the United States for arbitration. Negotiations between the Ecuadorian and Peruvian representatives began in Washington on September 30, 1935. These negotiations were long and tiresome. Both sides logically presented their cases, but no one seemed to give up their claims. Then on February 6, 1937, Ecuador presented a transactional line which Peru rejected the next day. The negotiations turned into intense arguments during the next 7 months and finally on September 29, 1937, the Peruvian representatives decided to break off the negotiations without submitting the dispute to arbitration because the direct negotiations were going nowhere.
+
+Four years later in 1941, amid fast-growing tensions within disputed territories around the Zarumilla River, war broke out with Peru. Peru claimed that Ecuador's military presence in Peruvian-claimed territory was an invasion; Ecuador, for its part, claimed that Peru had recently invaded Ecuador around the Zarumilla River and that Peru since Ecuador's independence from Spain has systematically occupied Tumbez, Jaen, and most of the disputed territories in the Amazonian Basin between the Putomayo and Marañon Rivers. In July 1941, troops were mobilized in both countries. Peru had an army of 11,681 troops who faced a poorly supplied and inadequately armed Ecuadorian force of 2,300, of which only 1,300 were deployed in the southern provinces. Hostilities erupted on July 5, 1941, when Peruvian forces crossed the Zarumilla river at several locations, testing the strength and resolve of the Ecuadorian border troops. Finally, on July 23, 1941, the Peruvians launched a major invasion, crossing the Zarumilla river in force and advancing into the Ecuadorian province of El Oro.
+
+During the course of the Ecuadorian–Peruvian War, Peru gained control over part of the disputed territory and some parts of the province of El Oro, and some parts of the province of Loja, demanding that the Ecuadorian government give up its territorial claims. The Peruvian Navy blocked the port of Guayaquil, almost cutting all supplies to the Ecuadorian troops. After a few weeks of war and under pressure by the United States and several Latin American nations, all fighting came to a stop. Ecuador and Peru came to an accord formalized in the Rio Protocol, signed on January 29, 1942, in favor of hemispheric unity against the Axis Powers in World War II favouring Peru with the territory they occupied at the time the war came to an end.
+
+The 1944 Glorious May Revolution followed a military-civilian rebellion and a subsequent civic strike which successfully removed Carlos Arroyo del Río as a dictator from Ecuador's government. However, a post-Second World War recession and popular unrest led to a return to populist politics and domestic military interventions in the 1960s, while foreign companies developed oil resources in the Ecuadorian Amazon. In 1972, construction of the Andean pipeline was completed. The pipeline brought oil from the east side of the Andes to the coast, making Ecuador South America's second largest oil exporter. The pipeline in southern Ecuador did nothing to resolve tensions between Ecuador and Peru, however.
+
+The Rio Protocol failed to precisely resolve the border along a little river in the remote Cordillera del Cóndor region in southern Ecuador. This caused a long-simmering dispute between Ecuador and Peru, which ultimately led to fighting between the two countries; first a border skirmish in January–February 1981 known as the Paquisha Incident, and ultimately full-scale warfare in January 1995 where the Ecuadorian military shot down Peruvian aircraft and helicopters and Peruvian infantry marched into southern Ecuador. Each country blamed the other for the onset of hostilities, known as the Cenepa War. Sixto Durán Ballén, the Ecuadorian president, famously declared that he would not give up a single centimeter of Ecuador. Popular sentiment in Ecuador became strongly nationalistic against Peru: graffiti could be seen on the walls of Quito referring to Peru as the "Cain de Latinoamérica", a reference to the murder of Abel by his brother Cain in the Book of Genesis.[27]
+
+Ecuador and Peru signed the Brasilia Presidential Act peace agreement on October 26, 1998, which ended hostilities, and effectively put an end to the Western Hemisphere's longest running territorial dispute.[28] The Guarantors of the Rio Protocol (Argentina, Brazil, Chile, and the United States of America) ruled that the border of the undelineated zone was to be set at the line of the Cordillera del Cóndor. While Ecuador had to give up its decades-old territorial claims to the eastern slopes of the Cordillera, as well as to the entire western area of Cenepa headwaters, Peru was compelled to give to Ecuador, in perpetual lease but without sovereignty, 1 km2 (0.39 sq mi) of its territory, in the area where the Ecuadorian base of Tiwinza – focal point of the war – had been located within Peruvian soil and which the Ecuadorian Army held during the conflict. The final border demarcation came into effect on May 13, 1999 and the multi-national MOMEP (Military Observer Mission for Ecuador and Peru) troop deployment withdrew on June 17, 1999.[28]
+
+In 1972, a "revolutionary and nationalist" military junta overthrew the government of Velasco Ibarra. The coup d'état was led by General Guillermo Rodríguez and executed by navy commander Jorge Queirolo G. The new president exiled José María Velasco to Argentina. He remained in power until 1976, when he was removed by another military government. That military junta was led by Admiral Alfredo Poveda, who was declared chairman of the Supreme Council. The Supreme Council included two other members: General Guillermo Durán Arcentales and General Luis Leoro Franco. The civil society more and more insistently called for democratic elections. Colonel Richelieu Levoyer, Government Minister, proposed and implemented a Plan to return to the constitutional system through universal elections. This plan enabled the new democratically elected president to assume the duties of the executive office.
+
+Elections were held on April 29, 1979, under a new constitution. Jaime Roldós Aguilera was elected president, garnering over one million votes, the most in Ecuadorian history. He took office on August 10, as the first constitutionally elected president after nearly a decade of civilian and military dictatorships. In 1980, he founded the Partido Pueblo, Cambio y Democracia (People, Change, and Democracy Party) after withdrawing from the Concentración de Fuerzas Populares (Popular Forces Concentration) and governed until May 24, 1981, when he died along with his wife and the minister of defense, Marco Subia Martinez, when his Air Force plane crashed in heavy rain near the Peruvian border. Many people believe that he was assassinated by the CIA,[citation needed] given the multiple death threats leveled against him because of his reformist agenda, deaths in automobile crashes of two key witnesses before they could testify during the investigation, and the sometimes contradictory accounts of the incident.
+
+Roldos was immediately succeeded by Vice President Osvaldo Hurtado, who was followed in 1984 by León Febres Cordero from the Social Christian Party. Rodrigo Borja Cevallos of the Democratic Left (Izquierda Democrática, or ID) party won the presidency in 1988, running in the runoff election against Abdalá Bucaram (brother in law of Jaime Roldos and founder of the Ecuadorian Roldosist Party). His government was committed to improving human rights protection and carried out some reforms, notably an opening of Ecuador to foreign trade. The Borja government concluded an accord leading to the disbanding of the small terrorist group, "¡Alfaro Vive, Carajo!" ("Alfaro Lives, Dammit!"), named after Eloy Alfaro. However, continuing economic problems undermined the popularity of the ID, and opposition parties gained control of Congress in 1999.
+
+The emergence of the Amerindian population as an active constituency has added to the democratic volatility of the country in recent years. The population has been motivated by government failures to deliver on promises of land reform, lower unemployment and provision of social services, and historical exploitation by the land-holding elite. Their movement, along with the continuing destabilizing efforts by both the elite and leftist movements, has led to a deterioration of the executive office. The populace and the other branches of government give the president very little political capital, as illustrated by the most recent removal of President Lucio Gutiérrez from office by Congress in April 2005. Vice President Alfredo Palacio took his place and remained in office until the presidential election of 2006, in which Rafael Correa gained the presidency.[29]
+
+In December 2008, president Correa declared Ecuador's national debt illegitimate, based on the argument that it was odious debt contracted by corrupt and despotic prior regimes. He announced that the country would default on over $3 billion worth of bonds; he then pledged to fight creditors in international courts and succeeded in reducing the price of outstanding bonds by more than 60%.[30] He brought Ecuador into the Bolivarian Alliance for the Americas in June 2009. To date, Correa's administration has succeeded in reducing the high levels of poverty and unemployment in Ecuador.[31][32][33][34][35]
+
+After being elected in 2017, President Lenin Moreno's government adopted economically liberal policies: reduction of public spending, trade liberalization, flexibility of the labour code, etc. He also left the left-wing Bolivarian Alliance for the Americas in August 2018. The Productive Development Act enshrines an austerity policy, and reduces the development and redistribution policies of the previous mandate. In the area of taxes, the authorities aim to "encourage the return of investors" by granting amnesty to fraudsters and proposing measures to reduce tax rates for large companies. In addition, the government waives the right to tax increases in raw material prices and foreign exchange repatriations.[36]
+
+A  series of protests began on 3 October 2019 against the end of fuel subsidies and austerity measures adopted by President of Ecuador Lenín Moreno and his administration. On 10 October, protesters overran the capital Quito causing the Government of Ecuador to relocate to Guayaquil,[37] but it was reported that the government still had plans to return to Quito.[38]
+
+The Ecuadorian State consists of five branches of government: the Executive Branch, the Legislative Branch, the Judicial Branch, the Electoral Branch, and Transparency and Social Control.
+
+Ecuador is governed by a democratically elected president, for a four-year term. The current president of Ecuador, Lenín Moreno, exercises his power from the presidential Palacio de Carondelet in Quito. The current constitution was written by the Ecuadorian Constituent Assembly elected in 2007, and was approved by referendum in 2008. Since 1936, voting is compulsory for all literate persons aged 18–65, optional for all other citizens.[39]
+
+The executive branch includes 23 ministries. Provincial governors and councilors (mayors, aldermen, and parish boards) are directly elected. The National Assembly of Ecuador meets throughout the year except for recesses in July and December. There are thirteen permanent committees. Members of the National Court of Justice are appointed by the National Judicial Council for nine-year terms.
+
+The executive branch is led by the president, an office currently held by Lenín Moreno. He is accompanied by the vice-president, currently Otto Sonnenholzner, elected for four years (with the ability to be re-elected only once). As head of state and chief government official, he is responsible for public administration including the appointing of national coordinators, ministers, ministers of State and public servants. The executive branch defines foreign policy, appoints the Chancellor of the Republic, as well as ambassadors and consuls, being the ultimate authority over the Armed Forces of Ecuador, National Police of Ecuador, and appointing authorities. The acting president's wife receives the title of First Lady of Ecuador.
+
+The legislative branch is embodied by the National Assembly, which is headquartered in the city of Quito in the Legislative Palace, and consists of 137 assemblymen, divided into ten committees and elected for a four-year term. Fifteen national constituency elected assembly, two Assembly members elected from each province and one for every 100,000 inhabitants or fraction exceeding 150,000, according to the latest national population census. In addition, statute determines the election of assembly of regions and metropolitan districts.
+
+Ecuador's judiciary has as its main body the Judicial Council, and also includes the National Court of Justice, provincial courts, and lower courts. Legal representation is made by the Judicial Council.
+The National Court of Justice is composed of 21 judges elected for a term of nine years. Judges are renewed by thirds every three years pursuant to the Judicial Code. These are elected by the Judicial Council on the basis of opposition proceedings and merits.
+The justice system is buttressed by the independent offices of public prosecutor and the public defender. Auxiliary organs are as follows: notaries, court auctioneers, and court receivers. Also there is a special legal regime for Amerindians.
+
+The electoral system functions by authorities which enter only every four years or when elections or referendums occur. Its main functions are to organize, control elections, and punish the infringement of electoral rules. Its main body is the National Electoral Council, which is based in the city of Quito, and consists of seven members of the political parties most voted, enjoying complete financial and administrative autonomy. This body, along with the electoral court, forms the Electoral Branch which is one of Ecuador's five branches of government.
+
+The Transparency and Social Control consists of the Council of Citizen Participation and Social Control, an ombudsman, the Comptroller General of the State, and the superintendents. Branch members hold office for five years. This branch is responsible for promoting transparency and control plans publicly, as well as plans to design mechanisms to combat corruption, as also designate certain authorities, and be the regulatory mechanism of accountability in the country.
+
+UN's Human Rights Council's (HRC) Universal Periodic Review (UPR) has treated the restrictions on freedom of expression and efforts to control NGOs and recommended that Ecuador should stop the criminal sanctions for the expression of opinions, and delay in implementing judicial reforms. Ecuador rejected the recommendation on decriminalization of libel.[40]
+
+According to Human Rights Watch (HRW) President Correa has intimidated journalists and subjected them to "public denunciation and retaliatory litigation". The sentences to journalists have been years of imprisonment and millions of dollars of compensation, even though defendants have been pardoned.[40] Correa has stated he was only seeking a retraction for slanderous statements.[41]
+
+According to HRW, Correa's government has weakened the freedom of press and independence of the judicial system. In Ecuador's current judicial system, judges are selected in a contest of merits, rather than government appointments. However, the process of selection has been criticized as biased and subjective. In particular, the final interview is said to be given "excessive weighing". Judges and prosecutors that have made decisions in favor of Correa in his lawsuits have received permanent posts, while others with better assessment grades have been rejected.[40][42]
+
+The laws also forbid articles and media messages that could favor or disfavor some political message or candidate. In the first half of 2012, twenty private TV or radio stations were closed down.[40]
+
+In July 2012, the officials warned the judges that they would be sanctioned and possibly dismissed if they allowed the citizens to appeal to the protection of their constitutional rights against the state.[40]
+
+People engaging in public protests against environmental and other issues are prosecuted for "terrorism and sabotage", which may lead to an eight-year prison sentence.[40]
+
+Ecuador joined the Organization of Petroleum Exporting Countries (OPEC) in 1973 and suspended its membership in 1992. Under President Rafael Correa, the country returned to OPEC before leaving again in 2020 under the instruction of President Moreno, citing its desire to increase crude oil importation to gain more revenue.[43][44]
+
+In Antarctica, Ecuador has maintained a peaceful research station for scientific study as a member nation of the Antarctica Treaty. Ecuador has often placed great emphasis on multilateral approaches to international issues. Ecuador is a member of the United Nations (and most of its specialized agencies) and a member of many regional groups, including the Rio Group, the Latin American Economic System, the Latin American Energy Organization, the Latin American Integration Association, the Andean Community of Nations, and the Bank of the South (Spanish: Banco del Sur or BancoSur).
+
+In 2017, the Ecuadorian parliament adopted a Law on human mobility.[45]
+
+The International Organization for Migration lauds Ecuador as the first state to have established the promotion of the concept of universal citizenship in its constitution, aiming to promote the universal recognition and protection of the human rights of migrants.[46] In 2017, Ecuador signed the UN treaty on the Prohibition of Nuclear Weapons.[47]
+
+Ecuador is divided into 24 provinces (Spanish: provincias), each with its own administrative capital:
+
+The provinces are divided into cantons and further subdivided into parishes (parroquias).
+
+Regionalization, or zoning, is the union of two or more adjoining provinces in order to decentralize the administrative functions of the capital, Quito.
+In Ecuador, there are seven regions, or zones, each shaped by the following provinces:
+
+Quito and Guayaquil are Metropolitan Districts. Galápagos, despite being included within Region 5,[49] is also under a special unit.[50]
+
+The Ecuadorian Armed Forces (Fuerzas Armadas de la Republica de Ecuador), consists of the Army, Air Force, and Navy and have the stated responsibility for the preservation of the integrity and national sovereignty of the national territory.
+
+The military tradition starts in Gran Colombia, where a sizable army was stationed in Ecuador due to border disputes with Peru, which claimed territories under its political control when it was a Spanish vice-royalty. Once Gran Colombia was dissolved after the death of Simón Bolívar in 1830, Ecuador inherited the same border disputes and had the need of creating its own professional military force. So influential was the military in Ecuador in the early republican period that its first decade was under the control of General Juan José Flores, first president of Ecuador of Venezuelan origin. General Jose Ma. Urbina and General Robles are examples of military figures who became presidents of the country in the early republican period.
+
+Due to the continuous border disputes with Peru, finally settled in the early 2000s, and due to the ongoing problem with the Colombian guerrilla insurgency infiltrating Amazonian provinces, the Ecuadorian Armed Forces has gone through a series of changes. In 2009, the new administration at the Defense Ministry launched a deep restructuring within the forces, increasing spending budget to $1,691,776,803, an increase of 25%.[51]
+
+The icons of the Ecuadorian military forces are Marshall Antonio José de Sucre and General Eloy Alfaro.
+
+The Military Academy General Eloy Alfaro (c. 1838) located in Quito is in charge to graduate the army officers.[52]
+
+The Ecuadorian Navy Academy (c. 1837), located in Salinas graduates the navy officers.[53]
+
+The Air Academy "Cosme Rennella (c. 1920), also located in Salinas, graduates the air force officers.[54]
+
+Other training academies for different military specialties are found across the country.
+
+The IWIAS is a special force trained to perform exploration and militar activities. Is considered the best elite force of Ecuador and is conformed by indigenous of the Amazon who combine their inherital experience for jungle dominance with modern army tactics.
+
+Ecuador has a total area of 283,561 km2 (109,484 sq mi), including the Galápagos Islands. Of this, 276,841 km2 (106,889 sq mi) is land and 6,720 km2 (2,595 sq mi) water.[1] Ecuador is bigger than Uruguay, Suriname, Guyana and French Guyana in South America.
+
+Ecuador lies between latitudes 2°N and 5°S,
+bounded on the west by the Pacific Ocean, and has 2,337 km (1,452 mi) of coastline. It has 2,010 km (1,250 mi) of land boundaries, with Colombia in the north (with a 590 km (367 mi) border) and Peru in the east and south (with a 1,420 km (882 mi) border). It is the westernmost country that lies on the equator.[55]
+
+The country has four main geographic regions:
+
+Ecuador's capital is Quito, which is in the province of Pichincha in the Sierra region. Its largest city is Guayaquil, in the Guayas Province. Cotopaxi, just south of Quito, is one of the world's highest active volcanoes. The top of Mount Chimborazo (6,268 m, or 20,560 ft, above sea level), Ecuador's tallest mountain, is the most distant point from the center of the Earth on the Earth's surface because of the ellipsoid shape of the planet.[1]
+
+There is great variety in the climate, largely determined by altitude. It is mild year-round in the mountain valleys, with a humid subtropical climate in coastal areas and rainforest in lowlands. The Pacific coastal area has a tropical climate with a severe rainy season. The climate in the Andean highlands is temperate and relatively dry, and the Amazon basin on the eastern side of the mountains shares the climate of other rainforest zones.
+
+Because of its location at the equator, Ecuador experiences little variation in daylight hours during the course of a year. Both sunrise and sunset occur each day at the two six o'clock hours.[1]
+
+The Andes is the watershed divisor between the Amazon watershed, which runs to the east, and the Pacific, including the north–south rivers Mataje, Santiago, Esmeraldas, Chone, Guayas, Jubones, and Puyango-Tumbes.
+
+Almost all of the rivers in Ecuador form in the Sierra region and flow east toward the Amazon River or west toward the Pacific Ocean. The rivers rise from snowmelt at the edges of the snowcapped peaks or from the abundant precipitation that falls at higher elevations. In the Sierra region, the streams and rivers are narrow and flow rapidly over precipitous slopes. Rivers may slow and widen as they cross the hoyas yet become rapid again as they flow from the heights of the Andes to the lower elevations of the other regions. The highland rivers broaden as they enter the more level areas of the Costa and the Oriente.
+
+In the Costa, the external coast has mostly intermittent rivers that are fed by constant rains from December through May and become empty riverbeds during the dry season. The few exceptions are the longer, perennial rivers that flow throughout the external coast from the internal coast and La Sierra on their way to the Pacific Ocean. The internal coast, by contrast, is crossed by perennial rivers that may flood during the rainy season, sometimes forming swamps.
+
+Major rivers in the Oriente include the Pastaza, Napo, and Putumayo. The Pastaza is formed by the confluence of the Chambo and the Patate rivers, both of which rise in the Sierra. The Pastaza includes the Agoyan waterfall, which at sixty-one meters (200 feet) is the highest waterfall in Ecuador. The Napo rises near Mount Cotopaxi and is the major river used for transport in the eastern lowlands. The Napo ranges in width from 500 to 1,800 m (1,640 to 5,906 ft). In its upper reaches, the Napo flows rapidly until the confluence with one of its major tributaries, the Coca River, where it slows and levels off. The Putumayo forms part of the border with Colombia. All of these rivers flow into the Amazon River. The Galápagos Islands have no significant rivers. Several of the larger islands, however, have freshwater springs although they are surrounded by the Pacific Ocean.
+
+Ecuador is one of seventeen megadiverse countries in the world according to Conservation International,[19] and it has the most biodiversity per square kilometer of any nation.[56][57]
+
+Ecuador has 1,600 bird species (15% of the world's known bird species) in the continental area and 38 more endemic in the Galápagos. In addition to over 16,000 species of plants, the country has 106 endemic reptiles, 138 endemic amphibians, and 6,000 species of butterfly. The Galápagos Islands are well known as a region of distinct fauna, famous as the place of birth of Darwin's Theory of Evolution and a UNESCO World Heritage Site.[58]
+
+Ecuador has the first constitution to recognize the rights of nature.[59] The protection of the nation's biodiversity is an explicit national priority as stated in the National Plan of "Buen Vivir", or good living, Objective 4, "Guarantee the rights of nature", Policy 1: "Sustainably conserve and manage the natural heritage, including its land and marine biodiversity, which is considered a strategic sector".[56]  As of the writing of the Plan in 2008, 19% of Ecuador's land area was in a protected area; however, the Plan also states that 32% of the land must be protected in order to truly preserve the nation's biodiversity.[56] Current protected areas include 11 national parks, 10 wildlife refuges, 9 ecological reserves, and other areas.[60] A program begun in 2008, Sociobosque, is preserving another 2.3% of total land area (6,295 km2, or 629,500 ha) by paying private landowners or community landowners (such as Amerindian tribes) incentives to maintain their land as native ecosystems such as native forests or grasslands. Eligibility and subsidy rates for this program are determined based on the poverty in the region, the number of hectares that will be protected, and the type of ecosystem of the land to be protected, among other factors.[61]
+
+Despite being on the UNESCO list, the Galápagos are endangered by a range of negative environmental effects, threatening the existence of this exotic ecosystem.[62] Additionally, oil exploitation of the Amazon rainforest has led to the release of billions of gallons of untreated wastes, gas, and crude oil into the environment, contaminating ecosystems and causing detrimental health effects to Amerindian peoples.[63]
+
+Ecuador has a developing economy that is highly dependent on commodities, namely petroleum and agricultural products. The country is classified as an upper-middle-income country. Ecuador's economy is the eighth largest in Latin America and experienced an average growth of 4.6% between 2000 and 2006.[64][failed verification] From 2007 to 2012, Ecuador's GDP grew at an annual average of 4.3 percent, above the average for Latin America and the Caribbean, which was 3.5%, according to the United Nations' Economic Commission for Latin American and the Caribbean (ECLAC).[65] Ecuador was able to maintain relatively superior growth during the crisis. In January 2009, the Central Bank of Ecuador (BCE) put the 2010 growth forecast at 6.88%.[66] In 2011, its GDP grew at 8% and ranked 3rd highest in Latin America, behind Argentina (2nd) and Panama (1st).[67] Between 1999 and 2007, GDP doubled, reaching $65,490 million according to BCE.[68]
+The inflation rate until January 2008, was about 1.14%, the highest in the past year, according to the government.[69][70] The monthly unemployment rate remained at about 6 and 8 percent from December 2007 until September 2008; however, it went up to about 9 percent in October and dropped again in November 2008 to 8 percent.[71] Unemployment mean annual rate for 2009 in Ecuador was 8.5% because the global economic crisis continued to affect the Latin American economies. From this point, unemployment rates started a downward trend: 7.6% in 2010, 6.0% in 2011, and 4.8% in 2012.[72]
+
+The extreme poverty rate has declined significantly between 1999 and 2010.[73] In 2001, it was estimated at 40% of the population, while by 2011 the figure dropped to 17.4% of the total population.[74] This is explained to an extent by emigration and the economic stability achieved after adopting the U.S. dollar as official means of transaction (before 2000, the Ecuadorian sucre was prone to rampant inflation). However, starting in 2008, with the bad economic performance of the nations where most Ecuadorian emigrants work, the reduction of poverty has been realized through social spending, mainly in education and health.[75]
+
+Oil accounts for 40% of exports and contributes to maintaining a positive trade balance.[76] Since the late 1960s, the exploitation of oil increased production, and proven reserves are estimated at 6.51 billion barrels as of 2011[update].[77]
+
+The overall trade balance for August 2012 was a surplus of almost $390 million for the first six months of 2012, a huge figure compared with that of 2007, which reached only $5.7 million; the surplus had risen by about $425 million compared to 2006.[74] The oil trade balance positive had revenues of $3.295 million in 2008, while non-oil was negative, amounting to $2.842 million. The trade balance with the United States, Chile, the European Union, Bolivia, Peru, Brazil, and Mexico is positive. The trade balance with Argentina, Colombia, and Asia is negative.[78]
+
+In the agricultural sector, Ecuador is a major exporter of bananas (first place worldwide in production and export), flowers, and the seventh largest producer of cocoa.[79] Ecuador also produces coffee, rice, potatoes, cassava (manioc, tapioca), plantains and sugarcane; cattle, sheep, pigs, beef, pork and dairy products; fish, and shrimp; and balsa wood.[80] The country's vast resources include large amounts of timber across the country, like eucalyptus and mangroves.[81] Pines and cedars are planted in the region of La Sierra and walnuts, rosemary, and balsa wood in the Guayas River Basin.[82]
+The industry is concentrated mainly in Guayaquil, the largest industrial center, and in Quito, where in recent years the industry has grown considerably. This city is also the largest business center of the country.[83] Industrial production is directed primarily to the domestic market.[citation needed] Despite this, there is limited export of products produced or processed industrially.[citation needed] These include canned foods, liquor, jewelry, furniture, and more.[citation needed] A minor industrial activity is also concentrated in Cuenca.[84] Incomes from tourism has been increasing during the last few years because of the Government showing the variety of climates and the biodiversity of Ecuador.
+
+Ecuador has negotiated bilateral treaties with other countries, besides belonging to the Andean Community of Nations,[85] and an associate member of Mercosur.[86] It also serves on the World Trade Organization (WTO), in addition to the Inter-American Development Bank (IDB), World Bank, International Monetary Fund (IMF), Corporación Andina de Fomento (CAF) and other multilateral agencies.[87][88][89] In April 2007, Ecuador paid off its debt to the IMF, thus ending an era of interventionism of the Agency in the country.[90][91]
+The public finance of Ecuador consists of the Central Bank of Ecuador (BCE), the National Development Bank (BNF), the State Bank.
+
+The Ministry of Information and Tourism was created on August 10, 1992, at the beginning of the government of Sixto Durán Ballén, who viewed tourism as a fundamental activity for the economic and social development of the peoples. Faced with the growth of the tourism sector, in June 1994, the decision was taken to separate tourism from information, so that it is exclusively dedicated to promoting and strengthening this activity.
+
+Ecuador is a country with vast natural wealth. The diversity of its four regions has given rise to thousands of species of flora and fauna. It has around 1640 kinds of birds. The species of butterflies border the 4,500, the reptiles 345, the amphibians 358 and the mammals 258, among others. Not in vain, Ecuador is considered one of the 17 countries where the planet's highest biodiversity is concentrated, being also the largest country with diversity per km2 in the world. Most of its fauna and flora lives in 26 protected areas by the State. Also, it has a huge culture spectrum. Since 2007, with the government of Rafael Correa, the tourism brand "Ecuador Ama la Vida" has been transformed, with which the nation's tourism promotion would be sold. Focused on considering it as a country friendly and respectful of the nature, natural biodiversity and cultural diversity of the peoples. And for this, means of exploiting them are developed along with the private economy.
+
+The country has two cities UNESCO World Heritage Sites: Quito and Cuenca, as well as two natural UNESCO World Heritage Sites: the Galapagos Islands and Sangay National Park in addition to one World Biosphere Reserve, such as the Cajas Massif. Culturally, the Toquilla straw hat and the culture of the Zapara indigenous people are recognized. The most popular sites for national and foreign tourists have different nuances due to the various tourist activities offered by the country.
+
+Among the main tourist destinations are:
+
+The rehabilitation and reopening of the Ecuadorian railroad and use of it as a tourist attraction is one of the recent developments in transportation matters.[93]
+
+The roads of Ecuador in recent years have undergone important improvement. The major routes are Pan American (under enhancement from four to six lanes from Rumichaca to Ambato, the conclusion of 4 lanes on the entire stretch of Ambato and Riobamba and running via Riobamba to Loja). In the absence of the section between Loja and the border with Peru, there are the Route Espondilus and/or Ruta del Sol (oriented to travel along the Ecuadorian coastline) and the Amazon backbone (which crosses from north to south along the Ecuadorian Amazon, linking most and more major cities of it).
+
+Another major project is developing the road Manta – Tena, the highway Guayaquil – Salinas Highway Aloag Santo Domingo, Riobamba – Macas (which crosses Sangay National Park). Other new developments include the National Unity bridge complex in Guayaquil, the bridge over the Napo river in Francisco de Orellana, the Esmeraldas River Bridge in the city of the same name, and, perhaps the most remarkable of all, the Bahia – San Vincente Bridge, being the largest on the Latin American Pacific coast.
+
+Cuenca's tramway is the largest public transport system in the city and the first modern tramway in Ecuador. It was inaugurated on March 8, 2019. It has 20,4 km and 27 stations. It will transport 120 000 passagers daily. Its route starts in the south of Cuenca and ends in the north at the Parque Industrial neighbourhood.
+
+The Mariscal Sucre International Airport in Quito and the José Joaquín de Olmedo International Airport in Guayaquil have experienced a high increase in demand and have required modernization. In the case of Guayaquil it involved a new air terminal, once considered the best in South America and the best in Latin America[94] and in Quito where an entire new airport has been built in Tababela and was inaugurated in February 2013, with Canadian assistance. However, the main road leading from Quito city centre to the new airport will only be finished in late 2014, making current travelling from the airport to downtown Quito as long as two hours during rush hour.[95] Quito's old city-centre airport is being turned into parkland, with some light industrial use.
+
+Ecuador's population is ethnically diverse and the 2018 estimates put Ecuador's population at 17,084,358.[6][7] The largest ethnic group (as of 2010[update]) is the Mestizos, who are the descendants of Spanish colonists that interbred with Amerindian peoples, and constitute about 71% of the population. The White Ecuadorians (White Latin American) are a minority accounting for 6.1% of the population of Ecuador and can be found throughout all of Ecuador primarily around the urban areas. Even though Ecuador's white population during its colonial era were mainly descendants from Spain, today Ecuador's white population is a result of a mixture of European immigrants, predominantly from Spain with people from Italy, Germany, France, and Switzerland who have settled in the early 20th century. Ecuador also has people of middle eastern extraction that have also joined the ranks of the white minority. These include economically well off immigrants of Lebanese and Palestinian descent, who are either Christian or Muslim (Islam in Ecuador). In addition, there is a small European Jewish (Ecuadorian Jews) population, which is based mainly in Quito and to a lesser extent in Guayaquil.[3] Amerindians account for 7% of the current population. The mostly rural Montubio population of the coastal provinces of Ecuador, who might be classified as Pardo account for 7.4% of the population. The Afro-Ecuadorians are a minority population (7%) in Ecuador, that includes the Mulattos and zambos, and are largely based in the Esmeraldas province and to a lesser degree in the predominantly Mestizo provinces of Coastal Ecuador - Guayas and Manabi. In the Highland Andes where a predominantly Mestizo, white and Amerindian population exist, the African presence is almost non-existent except for a small community in the province of Imbabura called Chota Valley.
+
+According to the Ecuadorian National Institute of Statistics and Census, 91.95% of the country's population have a religion, 7.94% are atheists and 0.11% are agnostics. Among the people that have a religion, 80.44% are Roman Catholic Latin Rite (see List of Roman Catholic dioceses in Ecuador), 11.30% are Evangelical Protestants, 1.29% are Jehovah's Witnesses and 6.97% other (mainly Jewish, Buddhists and Latter-day Saints).[97][98]
+
+In the rural parts of Ecuador, Amerindian beliefs and Catholicism are sometimes syncretized. Most festivals and annual parades are based on religious celebrations, many incorporating a mixture of rites and icons.[citation needed]
+
+There is a small number of Eastern Orthodox Christians, Amerindian religions, Muslims (see Islam in Ecuador), Buddhists and Bahá'í. According to their own estimates, The Church of Jesus Christ of Latter-day Saints accounts for about 1.4% of the population, or 211,165 members at the end of 2012.[99] According to their own sources, in 2017 there were 92,752 Jehovah's Witnesses in the country.[100]
+
+The first Jews arrived in Ecuador in the 16th and 17th centuries. Most of them are Sephardic Anusim (Crypto-Jews) and many still speak Judaeo-Spanish (Ladino) language.[101][citation needed] Today the Jewish Community of Ecuador (Comunidad Judía del Ecuador) has its seat in Quito and has approximately 200 members. Nevertheless, this number is declining because young people leave the country for the United States or Israel. The Community has a Jewish Center with a synagogue, a country club, and a cemetery. It supports the "Albert Einstein School", where Jewish history, religion, and Hebrew classes are offered. There are very small communities in Cuenca. The "Comunidad de Culto Israelita" reunites the Jews of Guayaquil. This community works independently from the "Jewish Community of Ecuador" and is composed of only 30 people.[102]
+
+Iglesia de San Sebastián church in Cuenca
+
+Iglesia y Convento de San Francisco in Quito
+
+The Ecuadorian constitution recognizes the "pluri-nationality" of those who want to exercise their affiliation with their native ethnic groups. Thus, in addition to criollos, mestizos, and Afro-Ecuadorians, some people belong to the Amerindian nations scattered in a few places in the coast, Quechua Andean villages, and the Amazonian jungle.
+
+According to a 2015 genealogical DNA testing, the average Ecuadorian is estimated to be 52.96% Native American, 41.77% European, and 5.26% Sub-Saharan African overall.[104]
+
+The majority of Ecuadorians live in the central provinces, the Andes mountains, or along the Pacific coast. The tropical forest region to the east of the mountains (El Oriente) remains sparsely populated and contains only about 3% of the population. Birth rate is 2-1 for each death. Marriages are usually from 14 and above using parental consent. About 12.4% of the population is married in the ages 15–19. Divorce rates are moderate.
+
+The five largest cities in the country are Quito (2.78 million inhabitants), Guayaquil (2.72 million inhabitants), Cuenca (636,996 inhabitants), Santo Domingo (458,580 inhabitants), and Ambato (387,309 inhabitants). While the most populated metropolitan areas of the country are those of Guayaquil, Quito, Cuenca, Manabí Centro (Portoviejo-Manta) and Ambato.[105]
+
+A small East Asian Latino community, estimated at 2,500, mainly consists of those of Japanese and Chinese descent, whose ancestors arrived as miners, farmhands and fishermen in the late 19th century.[1]
+
+In the early years of World War II, Ecuador still admitted a certain number of immigrants, and in 1939, when several South American countries refused to accept 165 Jewish refugees from Germany aboard the ship Koenigstein, Ecuador granted them entry permits.[107]
+
+In recent years, Ecuador has grown in popularity among North American expatriates.[108] They're drawn there by the authentic cultural experience and beautiful natural surroundings. Also, Ecuador's favorable residency options make for an easy transition for those who decide to settle there indefinitely.
+
+Another perk that draws many expats to Ecuador is its low cost of living. Since everything from gas to groceries costs far less than in North America, it is a popular choice for those who are looking to make the most of their retirement budget.[109]
+
+Even real estate in Ecuador is much less than its tropical counterparts. However, as more and more North Americans are discovering Ecuador's potential, property prices are beginning to rise from where they were a decade ago, particularly in the areas that are popular among expats and tourists.
+
+Ecuador's mainstream culture is defined by its Hispanic mestizo majority, and, like their ancestry, it is traditionally of Spanish heritage, influenced in different degrees by Amerindian traditions and in some cases by African elements. The first and most substantial wave of modern immigration to Ecuador consisted of Spanish colonists, following the arrival of Europeans in 1499. A lower number of other Europeans and North Americans migrated to the country in the late 19th and early twentieth centuries and, in smaller numbers, Poles, Lithuanians, English, Irish, and Croats during and after the Second World War.
+
+Since African slavery was not the workforce of the Spanish colonies in the Andes Mountains, given the subjugation of the Amerindian people through proselytization and encomiendas, the minority population of African descent is mostly found in the coastal northern province of Esmeraldas. This is largely owing to the 17th-century shipwreck of a slave-trading galleon off the northern coast of Ecuador. The few black African survivors swam to the shore and penetrated the then-thick jungle under the leadership of Anton, the chief of the group, where they remained as free men maintaining their original culture, not influenced by the typical elements found in other provinces of the coast or in the Andean region. A little later, freed slaves from Colombia known as cimarrones joined them. In the small Chota Valley of the province of Imbabura exists a small community of Africans among the province's predominantly mestizo population. These blacks are descendants of Africans, who were brought over from Colombia by Jesuits to work their colonial sugar plantations as slaves. As a general rule, small elements of zambos and mulattoes coexisted among the overwhelming mestizo population of coastal Ecuador throughout its history as gold miners in Loja, Zaruma, and Zamora and as shipbuilders and plantation workers around the city of Guayaquil. Today you can find a small community of Africans in the Catamayo valley of the predominantly mestizo population of Loja.
+
+Ecuador's Amerindian communities are integrated into the mainstream culture to varying degrees,[110] but some may also practice their own native cultures, particularly the more remote Amerindian communities of the Amazon basin. Spanish is spoken as the first language by more than 90% of the population and as a first or second language by more than 98%. Part of Ecuador's population can speak Amerindian languages, in some cases as a second language. Two percent of the population speak only Amerindian languages.
+
+Most Ecuadorians speak Spanish as their first language, with its ubiquity permeating and dominating most of the country, though there are many who speak an Amerindian language, such as Kichwa (also spelled Quechua), which is one of the Quechuan languages and is spoken by approximately 2.5 million people in Ecuador, Bolivia, Colombia, and Peru.[111] Other Amerindian languages spoken in Ecuador include Awapit (spoken by the Awá), A'ingae (spoken by the Cofan), Shuar Chicham (spoken by the Shuar), Achuar-Shiwiar (spoken by the Achuar and the Shiwiar), Cha'palaachi (spoken by the Chachi), Tsa'fiki (spoken by the Tsáchila), Paicoca (spoken by the Siona and Secoya), and Wao Tededeo (spoken by the Waorani). Use of these Amerindian languages are, however, gradually diminishing due to Spanish's widespread use in education. Though most features of Ecuadorian Spanish are universal to the Spanish-speaking world, there are several idiosyncrasies.
+
+The music of Ecuador has a long history. Pasillo is a genre of indigenous Latin music. In Ecuador it is the "national genre of music". Through the years, many cultures have brought their influences together to create new types of music. There are also different kinds of traditional music like albazo, pasacalle, fox incaico, tonada, capishca, Bomba (highly established in Afro-Ecuadorian societies), and so on. Tecnocumbia and Rockola are clear examples of the influence of foreign cultures. One of the most traditional forms of dancing in Ecuador is Sanjuanito. It is originally from northern Ecuador (Otavalo-Imbabura). Sanjuanito is a type of dance music played during festivities by the mestizo and Amerindian communities. According to the Ecuadorian musicologist Segundo Luis Moreno, Sanjuanito was danced by Amerindian people during San Juan Bautista's birthday. This important date was established by the Spaniards on June 24, coincidentally the same date when Amerindian people celebrated their rituals of Inti Raymi.
+
+Ecuadorian cuisine is diverse, varying with the altitude and associated agricultural conditions. Most regions in Ecuador follow the traditional three-course meal of soup, a course that includes rice and a protein, and then dessert and coffee to finish. Supper is usually lighter and sometimes consists only of coffee or herbal tea with bread.[citation needed]
+
+In the highland region; grilled pork, chicken, beef, and cuy (guinea pig) are popular and are served with a variety of grains (especially rice and mote) or potatoes.[citation needed]
+
+In the coastal region, seafood is very popular, with fish, shrimp, and ceviche being key parts of the diet. Generally, ceviches are served with fried plantain (chifles or patacones), popcorn, or tostado. Plantain- and peanut-based dishes are the basis of most coastal meals. Encocados (dishes that contain a coconut sauce) are also very popular. Churrasco is a staple food of the coastal region, especially Guayaquil. Arroz con menestra y carne asada (rice with beans and grilled beef) is one of the traditional dishes of Guayaquil, as is fried plantain, which is often served with it. This region is a leading producer of bananas, cocoa beans (to make chocolate), shrimp, tilapia, mango, and passion fruit, among other products.[citation needed]
+
+In the Amazon region, a dietary staple is the yuca, elsewhere called cassava. Many fruits are available in this region, including bananas, tree grapes, and peach palms.[citation needed]
+
+Early literature in colonial Ecuador, as in the rest of Spanish America, was influenced by the Spanish Golden Age. One of the earliest examples is Jacinto Collahuazo,[112] an Amerindian chief of a northern village in today's Ibarra, born in the late 1600s. Despite the early repression and discrimination of the native people by the Spanish, Collahuazo learned to read and write in Castilian, but his work was written in Quechua. The use of Quipu was banned by the Spanish,[113] and in order to preserve their work, many Inca poets had to resort to the use of the Latin alphabet to write in their native Quechua language. The history behind the Inca drama "Ollantay", the oldest literary piece in existence for any Amerindian language in America,[114] shares some similarities with the work of Collahuazo. Collahuazo was imprisoned and all of his work burned. The existence of his literary work came to light many centuries later, when a crew of masons was restoring the walls of a colonial church in Quito and found a hidden manuscript. The salvaged fragment is a Spanish translation from Quechua of the "Elegy to the Dead of Atahualpa",[112] a poem written by Collahuazo, which describes the sadness and impotence of the Inca people of having lost their king Atahualpa.
+
+Other early Ecuadorian writers include the Jesuits Juan Bautista Aguirre, born in Daule in 1725, and Father Juan de Velasco, born in Riobamba in 1727. De Velasco wrote about the nations and chiefdoms that had existed in the Kingdom of Quito (today Ecuador) before the arrival of the Spanish. His historical accounts are nationalistic, featuring a romantic perspective of precolonial history.
+
+Famous authors from the late colonial and early republic period include Eugenio Espejo, a printer and main author of the first newspaper in Ecuadorian colonial times; Jose Joaquin de Olmedo (born in Guayaquil), famous for his ode to Simón Bolívar titled Victoria de Junin; Juan Montalvo, a prominent essayist and novelist; Juan Leon Mera, famous for his work "Cumanda" or "Tragedy among Savages" and the Ecuadorian National Anthem; Juan A. Martinez with A la Costa';, Dolores Veintimilla;[115] and others.
+
+Contemporary Ecuadorian writers include the novelist Jorge Enrique Adoum; the poet Jorge Carrera Andrade; the essayist Benjamín Carrión; the poets Medardo Angel Silva, Jorge Carrera Andrade, and Luis Alberto Costales; the novelist Enrique Gil Gilbert; the novelist Jorge Icaza (author of the novel Huasipungo, translated to many languages); the short story author Pablo Palacio; and the novelist Alicia Yanez Cossio.
+
+In spite of Ecuador's considerable mystique, it is rarely featured as a setting in contemporary western literature. One exception is "The Ecuadorian Deception," a murder mystery/thriller authored by American Bear Mills. In it, George d'Hout, a website designer from the United States is lured under false pretenses to Guayaquil. A corrupt American archaeologist is behind the plot, believing d'Hout holds the keys to locating a treasure hidden by a buccaneer ancestor. The story is based on a real pirate by the name of George d'Hout who terrorized Guayaquil in the 16th Century.
+
+The best known art styles from Ecuador belonged to the Escuela Quiteña (Quito School), which developed from the 16th to 18th centuries, examples of which are on display in various old churches in Quito. Ecuadorian painters include Eduardo Kingman, Oswaldo Guayasamín, and Camilo Egas from the Indiginist Movement; Manuel Rendon, Jaime Zapata, Enrique Tábara, Aníbal Villacís, Theo Constanté, Luis Molinari, Araceli Gilbert, Judith Gutierrez, Felix Arauz, and Estuardo Maldonado from the Informalist Movement; Teddy Cobeña from expressionism and figurative style[116][117][118] and Luis Burgos Flor with his abstract, futuristic style. The Amerindian people of Tigua, Ecuador, are also world-renowned[citation needed] for their traditional paintings.
+
+Centro de Arte Contemporáneo, Quito
+
+Centro Cultural Metropolitano in the historic center of Quito
+
+The most popular sport in Ecuador, as in most South American countries, is football. Its best known professional teams include; Emelec from Guayaquil also the most popular team in Ecuador, Liga De Quito from Quito; Barcelona S.C. from Guayaquil;  Deportivo Quito, and El Nacional from Quito; Olmedo from Riobamba; and Deportivo Cuenca from Cuenca. Currently the most successful football team in Ecuador is LDU Quito, and it is the only Ecuadorian team that has won the Copa Libertadores, the Copa Sudamericana, and the Recopa Sudamericana; they were also runners-up in the 2008 FIFA Club World Cup. The matches of the Ecuadorian national team are the most-watched sporting events in the country.[citation needed] Ecuador has qualified for the final rounds of the 2002, the 2006, & the 2014 FIFA World Cups. The 2002 FIFA World Cup qualifying campaign was considered a huge success for the country and its inhabitants.[citation needed] The unusually high elevation of the home stadium in Quito often affects the performance of visiting teams. Ecuador finished in 2nd place in the CONMEBOL qualifiers behind Argentina and above the team that would become World Champions, Brazil. In the 2006 FIFA World Cup, Ecuador finished ahead of Poland and Costa Rica finishing second behind Germany in Group A in the 2006 World Cup. They were defeated by England in the second round.
+
+Ecuador has won two medals in the Olympic Games, both gained by 20-km (12 mi) racewalker Jefferson Pérez, who took gold in the 1996 games and silver 12 years later. Pérez also set a world best in the 2003 World Championships of 1:17:21 for the 20-km (12 mi) distance.[119]
+
+In world class professional cycling, Richard Carapaz became the first Ecuadorian to win a Grand Tour. He won the 2019 Giro d'Italia[120]
+
+The current structure of the Ecuadorian public health care system dates back to 1967.[121][122] The Ministry of the Public Health (Ministerio de Salud Pública del Ecuador) is the responsible entity of the regulation and creation of the public health policies and health care plans. The Minister of Public Health is appointed directly by the President of the Republic. The current minister, or Ecuadorian general surgeon, is Margarita Guevara.
+
+The philosophy of the Ministry of Public Health is the social support and service to the most vulnerable population,[123] and its main plan of action lies around communitarian health and preventive medicine.[123]
+
+The public healthcare system allows patients to be treated without an appointment in public general hospitals by general practitioners and specialists in the outpatient clinic (Consulta Externa) at no cost. This is done in the four basic specialties of pediatric, gynecology, clinic medicine, and surgery.[124] There are also public hospitals specialized to treat chronic diseases, target a particular group of the population, or provide better treatment in some medical specialties. Some examples in this group are the Gynecologic Hospitals, or Maternities, Children Hospitals, Geriatric Hospitals, and Oncology Institutes.
+
+Although well-equipped general hospitals are found in the major cities or capitals of provinces, there are basic hospitals in the smaller towns and canton cities for family care consultation and treatments in pediatrics, gynecology, clinical medicine, and surgery.[124]
+
+Community health care centers (Centros de Salud) are found inside metropolitan areas of cities and in rural areas. These are day hospitals that provide treatment to patients whose hospitalization is under 24 hours.[124]
+The doctors assigned to rural communities, where the Amerindian population can be substantial, have small clinics under their responsibility for the treatment of patients in the same fashion as the day hospitals in the major cities. The treatment in this case respects the culture of the community.[124]
+
+The public healthcare system should not be confused with the Ecuadorian Social Security healthcare service, which is dedicated to individuals with formal employment and who are affiliated obligatorily through their employers. Citizens with no formal employment may still contribute to the social security system voluntarily and have access to the medical services rendered by the social security system. The Ecuadorian Institute of Social Security (IESS) has several major hospitals and medical sub-centers under its administration across the nation.[125]
+
+Ecuador currently ranks 20, in most efficient health care countries, compared to 111 back in the year 2000.[126]  Ecuadorians have a life expectancy of 77.1 years.[127] The infant mortality rate is 13 per 1,000 live births,[128] a major improvement from approximately 76 in the early 1980s and 140 in 1950.[129] 23% of children under five are chronically malnourished.[128] Population in some rural areas have no access to potable water, and its supply is provided by mean of water tankers. There are 686 malaria cases per 100,000 people.[130] Basic health care, including doctor's visits, basic surgeries, and basic medications, has been provided free since 2008.[128] However, some public hospitals are in poor condition and often lack necessary supplies to attend the high demand of patients. Private hospitals and clinics are well equipped but still expensive for the majority of the population.
+
+Between 2008 and 2016, new public hospitals have been built, the number of civil servants has increased significantly and salaries have been increased. In 2008, the government introduced universal and compulsory social security coverage. In 2015, corruption remains a problem. Overbilling is recorded in 20% of public establishments and in 80% of private establishments.[131]
+
+The Ecuadorian Constitution requires that all children attend school until they achieve a "basic level of education", which is estimated at nine school years.[133] In 1996, the net primary enrollment rate was 96.9%, and 71.8% of children stayed in school until the fifth grade / age 10.[133] The cost of primary and secondary education is borne by the government, but families often face significant additional expenses such as fees and transportation costs.[133]
+
+Provision of public schools falls far below the levels needed, and class sizes are often very large, and families of limited means often find it necessary to pay for education.[citation needed] In rural areas, only 10% of the children go on to high school.[134] The Ministry of Education states that the mean number of years completed is 6.7.[citation needed]
+
+Ecuador has 61 universities, many of which still confer terminal degrees according to the traditional Spanish education system,[135] honoring a long tradition of having some of the oldest universities in the Americas: University of San Fulgencio, founded in 1586 by the Augustines; San Gregorio Magno University, founded in 1651 by the Jesuits; and University of Santo Tomás of Aquino, founded in 1681 by the Dominican order.
+
+Among the traditional conferred terminal degrees can be noted the doctorate for medicine and law schools or engineering, physics, chemistry, or mathematics for polytechnic or technology institutes. These terminal degrees, as in the case of the PhD in other countries, were the main requirement for an individual to be accepted in academia as a professor or researcher. In the professional realm, a terminal degree granted by an accredited institution automatically provides a professional license to the individual.
+
+However, in 2004, the National Council of Higher Education (CONESUP), started the reorganization of all the degree-granting schemes of the accredited universities in order to pair them with foreign counterparts. The new structure of some careers caused the dropping of subjects, credits, or even the name of the previously conferred diplomas. The terminal degree in law, previously known as JD Juris Doctor (Doctor en Jurisprudencia) was replaced by the one of abogado (attorney) with the exception of the modification of the number of credits to equate it to an undergraduate degree. In the same fashion for medical school, the required time of education was considerably reduced from nine years (the minimum needed to obtain the title of MD in Medicine and Surgery) to almost five, with the provision that the diploma is not terminal anymore, and it is given with the title of médico (medic). Therefore, an MD or PhD in medicine is only to be obtained overseas until the universities adjust themselves to granting schemes and curriculum as in foreign counterparts. Nonetheless, a "médico" can start a career as family practitioner or general medicine physician.
+
+This new reorganization, although very ambitious, lacked the proper path to the homologation of diplomas for highly educated professionals graduated in the country or even for the ones graduated in foreign institutions. One of the points of conflict was the imposition of obtaining foreign degrees to current academicians. As today, a master's degree is a requirement to keep an academic position and at least a foreign PhD to attain or retain the status of rector (president of a university) or décano (dean). For Ecuadorian researchers and many academicians trained in the country, these regulations sounded illogical, disappointing, and unlawful since it appeared a question of a title name conflict rather than specialization or science advancement.
+
+A debate to modify this and other reforms, especially the one which granted control of the Higher Education System by the government, was practically passed with consensus by the multi-partisan National Assembly on August 4, 2010, but vetoed by President Rafael Correa, who wanted to keep the law strictly as it was originally redacted by his political party and SENPLADES (National Secretary of Planning and Development). Due to this change, there are many highly educated professionals and academicians under the old structure but estimated that only 87% of the faculty in public universities have already obtained a master's degree, and fewer than 5% have a PhD (although many of them already have Ecuadorian-granted doctorate degrees).
+
+About 300 institutes of higher education offer two to three years of post-secondary vocational or technical training.
+
+Ecuador is currently placed in 96th position of innovation in technology.[136] The most notable icons in Ecuadorian sciences are the mathematician and cartographer Pedro Vicente Maldonado, born in Riobamba in 1707, and the printer, independence precursor, and medical pioneer Eugenio Espejo, born in 1747 in Quito. Among other notable Ecuadorian scientists and engineers are Lieutenant Jose Rodriguez Labandera,[137] a pioneer who built the first submarine in Latin America in 1837; Reinaldo Espinosa Aguilar (1898–1950), a botanist and biologist of Andean flora; and José Aurelio Dueñas (1880–1961), a chemist and inventor of a method of textile serigraphy.
+
+The major areas of scientific research in Ecuador have been in the medical fields, tropical and infectious diseases treatments, agricultural engineering, pharmaceutical research, and bioengineering. Being a small country and a consumer of foreign technology, Ecuador has favored research supported by entrepreneurship in information technology. The antivirus program Checkprogram, banking protection system MdLock, and Core Banking Software Cobis are products of Ecuadorian development.[138]
+
+The scientific production in hard sciences has been limited due to lack of funding but focused around physics, statistics, and partial differential equations in mathematics.[citation needed] In the case of engineering fields, the majority of scientific production comes from the top three polytechnic institutions: Escuela Superior Politécnica del Litoral - ESPOL, Universidad de Las Fuerzas Armadas - ESPE, and Escuela Politécnica Nacional EPN. The Center for Research and Technology Development in Ecuador is an autonomous center for research and technology development funded by Senecyt.
+
+However, according to Nature, the multidisciplinary scientific journal the top 10 institutions that carry the most outstanding scientific contributions are: Yachay Tech University (Yachay Tech), Escuela Politécnica Nacional (EPN), and Universidad San Francisco de Quito (USFQ).[139]
+
+EPN is known for research and education in the  applied science, astronomy, atmospheric physics, engineering and physical sciences. The Geophysics Institute [140] monitors over the country's volcanoes in the Andes Mountains of Ecuador and in the Galápagos Islands, all of which is part of the Ring of Fire. EPN adopted the polytechnic university model that stresses laboratory instruction in applied science and engineering.
+
+The oldest observatory in South America is the Quito Astronomical Observatory and is located in Quito, Ecuador. The Quito Astronomical Observatory, which gives the global community of a Virtual Telescope System that is connected via the Internet and allows the world to watch by streaming, is managed by EPN.
+
+Contemporary Ecuadorian scientists who have been recognized by international institutions are Eugenia del Pino (born 1945), the first Ecuadorian to be elected to the United States National Academy of Science, and Arturo Villavicencio, who was part of the working group of the IPCC, which shared the 2007 Nobel Peace Prize with Al Gore for their dissemination of the effects of climate change.
+
+Currently, the politics of research and investigation are managed by the National Secretary of Higher Education, Science, and Technology (Senescyt).[141]

en/1794.html.txt

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+The equator of a rotating spheroid (such as a planet) is the parallel (circle of latitude) at which latitude is defined to be 0°. It is the imaginary line on the spheroid, equidistant from its poles, dividing it into northern and southern hemispheres. In other words, it is the intersection of the spheroid with the plane perpendicular to its axis of rotation and midway between its geographical poles.
+
+On Earth, the equator is about 40,075 km (24,901 mi) long, of which 78.8% lies across water and 21.3% over land. Indonesia is the country straddling the greatest length of the equatorial line across both land and sea.
+
+The name is derived from medieval Latin word aequator, in the phrase circulus aequator diei et noctis, meaning 'circle equalizing day and night', from the Latin word aequare meaning 'make equal'.[1]
+
+The latitude of the Earth's equator is, by definition, 0° (zero degrees) of arc. The equator is one of the five notable circles of latitude on Earth; the other four are both polar circles (the Arctic Circle and the Antarctic Circle) and both tropical circles (the Tropic of Cancer and the Tropic of Capricorn). The equator is the only line of latitude which is also a great circle—that is, one whose plane passes through the center of the globe. The plane of Earth's equator, when projected outwards to the celestial sphere, defines the celestial equator.
+
+In the cycle of Earth's seasons, the equatorial plane runs through the Sun twice per year: on the equinoxes in March and September. To a person on Earth, the Sun appears to travel above the equator (or along the celestial equator) at these times. Light rays from the Sun's center are perpendicular to Earth's surface at the point of solar noon on the equator.
+
+Locations on the equator experience the shortest sunrises and sunsets because the Sun's daily path is nearly perpendicular to the horizon for most of the year. The length of daylight (sunrise to sunset) is almost constant throughout the year; it is about 14 minutes longer than nighttime due to atmospheric refraction and the fact that sunrise begins (or sunset ends) as the upper limb, not the center, of the Sun's disk contacts the horizon.
+
+Earth bulges slightly at the equator; the "average" diameter of Earth is 12,750 km (7,920 mi), but the diameter at the equator is about 43 km (27 mi) greater than at the poles.[2]
+
+Sites near the equator, such as the Guiana Space Centre in Kourou, French Guiana, are good locations for spaceports as they have a fastest rotational speed of any latitude, 460 m/s. The added velocity reduces the fuel needed to launch spacecraft eastward (in the direction of Earth's rotation) to orbit, while simultaneously avoiding costly maneuvers to flatten inclination during missions such as the Apollo moon landings.[3]
+
+The precise location of the equator is not truly fixed; the true equatorial plane is perpendicular to the Earth's spin axis, which drifts about 9 metres (30 ft) during a year. This effect must be accounted for in detailed geophysical measurements.[citation needed]
+
+Geological samples show the equator significantly changed positions between 12 to 48 million years ago, as sediment deposited by ocean thermal currents at the equator have shifted. The deposits by thermal currents are determined by the axis of the earth, which determines solar coverage of the Earth’s surface. Changes in Earth axis can also be observed in the geographic layout of volcanic island chains, which are created by shifting hot spots under the Earth’s crust as the axis and crust move.[4]
+
+The International Association of Geodesy (IAG) and the International Astronomical Union (IAU) have chosen to use an equatorial radius of 6,378.1366 kilometres (3,963.1903 mi) (codified as the IAU 2009 value).[5] This equatorial radius is also in the 2003 and 2010 IERS Conventions.[6] It is also the equatorial radius used for the IERS 2003 ellipsoid.  If it were really circular, the length of the equator would then be exactly 2π times the radius, namely 40,075.0142 kilometres (24,901.4594 mi). The GRS 80 (Geodetic Reference System 1980) as approved and adopted by the IUGG at its Canberra, Australia meeting of 1979 has an equatorial radius of 6,378.137 kilometres (3,963.191 mi). The WGS 84 (World Geodetic System 1984) which is a standard for use in cartography, geodesy, and satellite navigation including GPS, also has an equatorial radius of 6,378.137 kilometres (3,963.191 mi). For both GRS 80 and WGS 84, this results in a length for the equator of 40,075.0167 km (24,901.4609 mi).
+
+The geographical mile is defined as one arc-minute of the equator, so it has different values depending on which radius is assumed. For example, by WSG-84, the distance is 1,855.3248 metres (6,087.024 ft), while by IAU-2000, it is 1,855.3257 metres (6,087.027 ft). This is a difference of less than one millimetre (0.039 in) over the total distance (approximately 1.86 kilometres or 1.16 miles).
+
+The earth is commonly modeled as a sphere flattened 0.336% along its axis. This makes the equator 0.16% longer than a meridian (a great circle passing through the two poles). The IUGG standard meridian is, to the nearest millimetre, 40,007.862917 kilometres (24,859.733480 mi), one arc-minute of which is 1,852.216 metres (6,076.82 ft), explaining the SI standardization of the nautical mile as 1,852 metres (6,076 ft), more than 3 metres (9.8 ft) less than the geographical mile.
+
+The sea-level surface of the Earth (the geoid) is irregular, so the actual length of the equator is not so easy to determine. Aviation Week and Space Technology on 9 October 1961 reported that measurements using the Transit IV-A satellite had shown the equatorial diameter from longitude 11° West to 169° East to be 1,000 feet (300 m) greater than its diameter ninety degrees away.[citation needed]
+
+The equator passes through the land of 11 countries. Starting at the Prime Meridian and heading eastwards, the equator passes through:
+
+Despite its name, no part of Equatorial Guinea lies on the equator. However, its island of Annobón is 155 km (96 mi) south of the equator, and the rest of the country lies to the north.
+
+Seasons result from the tilt of the Earth's axis compared to the plane of its revolution around the Sun. Throughout the year the northern and southern hemispheres are alternately turned either toward or away from the sun depending on Earth's position in its orbit. The hemisphere turned toward the sun receives more sunlight and is in summer, while the other hemisphere receives less sun and is in winter (see solstice).
+
+At the equinoxes, the Earth's axis is perpendicular to the sun rather than tilted toward or away, meaning that day and night are both about 12 hours long across the whole of the Earth.
+
+Near the equator, this means the variation in strength of solar radiation is different relative to the time of year than it is at higher latitudes: Maximum solar radiation is received during the equinoxes, when a place at the equator is under the subsolar point at high noon, and the intermediate seasons of spring and autumn occur at higher latitudes, and the minimum occurs during both solstices, when either pole is tilted towards or away from the sun, resulting in either summer or winter in both hemispheres. This also results in a corresponding movement of the equator away from the subsolar point, which is then situated over or near the relevant tropic circle. Nevertheless, temperatures are high year round due to the earth's axial tilt of 23.5° not being enough to create a low minimum midday declination to sufficiently weaken the sun's rays even during the solstices.
+
+Near the equator there is little temperature change throughout the year, though there may be dramatic differences in rainfall and humidity. The terms summer, autumn, winter and spring do not generally apply. Lowlands around the equator generally have a tropical rainforest climate, also known as an equatorial climate, though cold ocean currents cause some regions to have tropical monsoon climates with a dry season in the middle of the year, and the Somali Current generated by the Asian monsoon due to continental heating via the high Tibetan Plateau causes Greater Somalia to have an arid climate despite its equatorial location.
+
+Average annual temperatures in equatorial lowlands are around 31 °C (88 °F) during the afternoon and 23 °C (73 °F) around sunrise. Rainfall is very high away from cold ocean current upwelling zones, from 2,500 to 3,500 mm (100 to 140 in) per year. There are about 200 rainy days per year and average annual sunshine hours are around 2,000. Despite high year-round sea level temperatures, some higher altitudes such as the Andes and Mount Kilimanjaro have glaciers. The highest point on the equator is at the elevation of 4,690 metres (15,387 ft), at 0°0′0″N 77°59′31″W / 0.00000°N 77.99194°W / 0.00000; -77.99194 (highest point on the equator), found on the southern slopes of Volcán Cayambe [summit 5,790 metres (18,996 ft)] in Ecuador. This is slightly above the snow line and is the only place on the equator where snow lies on the ground. At the equator, the snow line is around 1,000 metres (3,300 ft) lower than on Mount Everest and as much as 2,000 metres (6,600 ft) lower than the highest snow line in the world, near the Tropic of Capricorn on Llullaillaco.
+
+There is a widespread maritime tradition of holding ceremonies to mark a sailor's first crossing of the equator. In the past, these ceremonies have been notorious for their brutality, especially in naval practice.[citation needed] Milder line-crossing ceremonies, typically featuring King Neptune, are also held for passengers' entertainment on some civilian ocean liners and cruise ships.[citation needed]

en/1795.html.txt

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+Coordinates: 2°00′S 77°30′W / 2.000°S 77.500°W / -2.000; -77.500
+
+in South America (grey)
+
+Ecuador (/ˈɛkwədɔːr/ (listen) EK-wə-dor; Spanish pronunciation: [ekwaˈðoɾ] (listen); Quechua: Ikwayur; Shuar: Ecuador or Ekuatur.),[12][13] officially the Republic of Ecuador (Spanish: República del Ecuador,  which literally translates as "Republic of the Equator"; Quechua: Ikwadur Ripuwlika; Shuar: Ekuatur Nunka),[14][15] is a country in northwestern South America, bordered by Colombia on the north, Peru on the east and south, and the Pacific Ocean on the west. Ecuador also includes the Galápagos Islands in the Pacific, about 1,000 kilometres (621 mi) west of the mainland. The capital is Quito.[16][17]
+
+The territories of modern-day Ecuador were once home to a variety of Amerindian groups that were gradually incorporated into the Inca Empire during the 15th century. The territory was colonized by Spain during the 16th century, achieving independence in 1820 as part of Gran Colombia, from which it emerged as its own sovereign state in 1830. The legacy of both empires is reflected in Ecuador's ethnically diverse population, with most of its 17.1 million people being mestizos, followed by large minorities of European, Amerindian, and African descendants. Spanish is the official language and is spoken by a majority of the population, though 13 Amerindian languages are also recognized, including Quechua and Shuar.
+
+The sovereign state of Ecuador is a middle-income representative democratic republic and a developing country[18] that is highly dependent on commodities, namely petroleum and agricultural products. It is governed as a democratic presidential republic. One of 17 megadiverse countries in the world,[19][20] Ecuador hosts many endemic plants and animals, such as those of the Galápagos Islands. In recognition of its unique ecological heritage, the new constitution of 2008 is the first in the world to recognize legally enforceable Rights of Nature, or ecosystem rights.[21] It also has the fifth lowest homicide rate in the Americas.[22] Between 2006 and 2016, poverty decreased from 36.7% to 22.5% and annual per capita GDP growth was 1.5 percent (as compared to 0.6 percent over the prior two decades). At the same time, the country's Gini index of economic inequality decreased from 0.55 to 0.47.[23]
+
+Various peoples had settled in the area of future Ecuador before the arrival of the Incas. The archeological evidence suggests that the Paleo-Indians' first dispersal into the Americas occurred near the end of the last glacial period, around 16,500–13,000 years ago. The first Indians who reached Ecuador may have journeyed by land from North and Central America or by boat down the Pacific Ocean coastline. Much later migrations to Ecuador may have come via the Amazon tributaries, others descended from northern South America, and others ascended from the southern part of South America through the Andes. They developed different languages while emerging as unique ethnic groups.
+
+Even though their languages were unrelated, these groups developed similar groups of cultures, each based in different environments. The people of the coast developed a fishing, hunting, and gathering culture; the people of the highland Andes developed a sedentary agricultural way of life, and the people of the Amazon basin developed a nomadic hunting-and-gathering mode of existence.
+
+Over time these groups began to interact and intermingle with each other so that groups of families in one area became one community or tribe, with a similar language and culture. Many civilizations arose in Ecuador, such as the Valdivia Culture and Machalilla Culture on the coast, the Quitus (near present-day Quito), and the Cañari (near present-day Cuenca). Each civilisation developed its own distinctive architecture, pottery, and religious interests.
+
+In the highland Andes mountains, where life was more sedentary, groups of tribes cooperated and formed villages; thus the first nations based on agricultural resources and the domestication of animals formed. Eventually, through wars and marriage alliances of their leaders, a group of nations formed confederations. One region consolidated under a confederation called the Shyris, which exercised organized trading and bartering between the different regions. Its political and military power came under the rule of the Duchicela blood-line.
+
+When the Incas arrived, they found that these confederations were so developed that it took the Incas two generations of rulers—Topa Inca Yupanqui and Huayna Capac—to absorb them into the Inca Empire. The native confederations that gave them the most problems were deported to distant areas of Peru, Bolivia, and north Argentina. Similarly, a number of loyal Inca subjects from Peru and Bolivia were brought to Ecuador to prevent rebellion. Thus, the region of highland Ecuador became part of the Inca Empire in 1463 sharing the same language.
+
+In contrast, when the Incas made incursions into coastal Ecuador and the eastern Amazon jungles of Ecuador, they found both the environment and indigenous people more hostile. Moreover, when the Incas tried to subdue them, these indigenous people withdrew to the interior and resorted to guerrilla tactics. As a result, Inca expansion into the Amazon Basin and the Pacific coast of Ecuador was hampered. The indigenous people of the Amazon jungle and coastal Ecuador remained relatively autonomous until the Spanish soldiers and missionaries arrived in force. The Amazonian people and the Cayapas of Coastal Ecuador were the only groups to resist Inca and Spanish domination, maintaining their language and culture well into the 21st century.
+
+Before the arrival of the Spaniards, the Inca Empire was involved in a civil war. The untimely death of both the heir Ninan Cuchi and the Emperor Huayna Capac, from a European disease that spread into Ecuador, created a power vacuum between two factions. The northern faction headed by Atahualpa claimed that Huayna Capac gave a verbal decree before his death about how the empire should be divided. He gave the territories pertaining to present-day Ecuador and northern Peru to his favorite son Atahualpa, who was to rule from Quito; and he gave the rest to Huáscar, who was to rule from Cuzco. He willed that his heart be buried in Quito, his favorite city, and the rest of his body be buried with his ancestors in Cuzco.
+
+Huáscar did not recognize his father's will, since it did not follow Inca traditions of naming an Inca through the priests. Huáscar ordered Atahualpa to attend their father's burial in Cuzco and pay homage to him as the new Inca ruler. Atahualpa, with a large number of his father's veteran soldiers, decided to ignore Huáscar, and a civil war ensued. A number of bloody battles took place until finally Huáscar was captured. Atahualpa marched south to Cuzco and massacred the royal family associated with his brother.
+
+In 1532, a small band of Spaniards headed by Francisco Pizarro landed in Tumbez and marched over the Andes Mountains until they reached Cajamarca, where the new Inca Atahualpa was to hold an interview with them. Valverde, the priest, tried to convince Atahualpa that he should join the Catholic Church and declare himself a vassal of Spain. This infuriated Atahualpa so much that he threw the Bible to the ground. At this point the enraged Spaniards, with orders from Valverde, attacked and massacred unarmed escorts of the Inca and captured Atahualpa. Pizarro promised to release Atahualpa if he made good his promise of filling a room full of gold. But, after a mock trial, the Spaniards executed Atahualpa by strangulation.
+
+New infectious diseases such as smallpox, endemic to the Europeans, caused high fatalities among the Amerindian population during the first decades of Spanish rule, as they had no immunity. At the same time, the natives were forced into the encomienda labor system for the Spanish. In 1563, Quito became the seat of a real audiencia (administrative district) of Spain and part of the Viceroyalty of Peru and later the Viceroyalty of New Granada.
+
+The 1797 Riobamba earthquake, which caused up to 40,000 casualties, was studied by Alexander von Humboldt, when he visited the area in 1801–1802.[24]
+
+After nearly 300 years of Spanish rule, Quito was still a small city numbering 10,000 inhabitants. On August 10, 1809, the city's criollos called for independence from Spain (first among the peoples of Latin America). They were led by Juan Pío Montúfar, Quiroga, Salinas, and Bishop Cuero y Caicedo. Quito's nickname, "Luz de América" ("Light of America"), is based on its leading role in trying to secure an independent, local government. Although the new government lasted no more than two months, it had important repercussions and was an inspiration for the independence movement of the rest of Spanish America. August 10 is now celebrated as Independence Day, a national holiday.[25]
+
+On October 9, 1820, the Department of Guayaquil became the first territory in Ecuador to gain its independence from Spain, and it spawned most of the Ecuadorian coastal provinces, establishing itself as an independent state. Its inhabitants celebrated what is now Ecuador's official Independence Day on May 24, 1822. The rest of Ecuador gained its independence after Antonio José de Sucre defeated the Spanish Royalist forces at the Battle of Pichincha, near Quito. Following the battle, Ecuador joined Simón Bolívar's Republic of Gran Colombia, also including modern-day Colombia, Venezuela and Panama. In 1830, Ecuador separated from Gran Colombia and became an independent republic.
+
+The 19th century was marked by instability for Ecuador with a rapid succession of rulers. The first president of Ecuador was the Venezuelan-born Juan José Flores, who was ultimately deposed, followed by several authoritarian leaders, such as Vicente Rocafuerte; José Joaquín de Olmedo; José María Urbina; Diego Noboa; Pedro José de Arteta; Manuel de Ascásubi; and Flores's own son, Antonio Flores Jijón, among others. The conservative Gabriel Garcia Moreno unified the country in the 1860s with the support of the Roman Catholic Church. In the late 19th century, world demand for cocoa tied the economy to commodity exports and led to migrations from the highlands to the agricultural frontier on the coast.
+
+Ecuador abolished slavery and freed its black slaves in 1851.[26]
+
+The Liberal Revolution of 1895 under Eloy Alfaro reduced the power of the clergy and the conservative land owners. This liberal wing retained power until the military "Julian Revolution" of 1925. The 1930s and 1940s were marked by instability and emergence of populist politicians, such as five-time President José María Velasco Ibarra.
+
+Brasilia Presidential Act
+
+Since Ecuador's separation from Colombia on May 13, 1830, its first President, General Juan José Flores, laid claim to the territory that was called the Real Audiencia of Quito, also referred to as the Presidencia of Quito. He supported his claims with Spanish Royal decrees or Real Cedulas, that delineated the borders of Spain's former overseas colonies. In the case of Ecuador, Flores-based Ecuador's de jure claims on the following cedulas - Real Cedula of 1563, 1739, and 1740; with modifications in the Amazon Basin and Andes Mountains that were introduced through the Treaty of Guayaquil (1829) which Peru reluctantly signed, after the overwhelmingly outnumbered Gran Colombian force led by Antonio José de Sucre defeated President and General La Mar's Peruvian invasion force in the Battle of Tarqui. In addition, Ecuador's eastern border with the Portuguese colony of Brazil in the Amazon Basin was modified before the wars of Independence by the First Treaty of San Ildefonso (1777) between the Spanish Empire and the Portuguese Empire. Moreover, to add legitimacy to his claims, on February 16, 1840, Flores signed a treaty with Spain, whereby Flores convinced Spain to officially recognize Ecuadorian independence and its sole rights to colonial titles over Spain's former colonial territory known anciently to Spain as the Kingdom and Presidency of Quito.
+
+Ecuador during its long and turbulent history has lost most of its contested territories to each of its more powerful neighbors, such as Colombia in 1832 and 1916, Brazil in 1904 through a series of peaceful treaties, and Peru after a short war in which the Protocol of Rio de Janeiro was signed in 1942.
+
+During the struggle for independence, before Peru or Ecuador became independent nations, a few areas of the former Vice Royalty of New Granada - Guayaquil, Tumbez, and Jaén - declared themselves independent from Spain. A few months later, a part of the Peruvian liberation army of San Martin decided to occupy the independent cities of Tumbez and Jaén with the intention of using these towns as springboards to occupy the independent city of Guayaquil and then to liberate the rest of the Audiencia de Quito (Ecuador). It was common knowledge among the top officers of the liberation army from the south that their leader San Martin wished to liberate present-day Ecuador and add it to the future republic of Peru, since it had been part of the Inca Empire before the Spaniards conquered it.
+
+However, Bolívar's intention was to form a new republic known as the Gran Colombia, out of the liberated Spanish territory of New Granada which consisted of Colombia, Venezuela, and Ecuador. San Martin's plans were thwarted when Bolívar, with the help of Marshal Antonio José de Sucre and the Gran Colombian liberation force, descended from the Andes mountains and occupied Guayaquil; they also annexed the newly liberated Audiencia de Quito to the Republic of Gran Colombia. This happened a few days before San Martin's Peruvian forces could arrive and occupy Guayaquil, with the intention of annexing Guayaquil to the rest of Audiencia of Quito (Ecuador) and to the future republic of Peru. Historic documents repeatedly stated that San Martin told Bolivar he came to Guayaquil to liberate the land of the Incas from Spain. Bolivar countered by sending a message from Guayaquil welcoming San Martin and his troops to Colombian soil.
+
+In the south, Ecuador had de jure claims to a small piece of land beside the Pacific Ocean known as Tumbes which lay between the Zarumilla and Tumbes rivers. In Ecuador's southern Andes Mountain region where the Marañon cuts across, Ecuador had de jure claims to an area it called Jaén de Bracamoros. These areas were included as part of the territory of Gran Colombia by Bolivar on December 17, 1819, during the Congress of Angostura when the Republic of Gran Colombia was created. Tumbes declared itself independent from Spain on January 17, 1821, and Jaen de Bracamoros on June 17, 1821, without any outside help from revolutionary armies. However, that same year, 1821, Peruvian forces participating in the Trujillo revolution occupied both Jaen and Tumbes. Some Peruvian generals, without any legal titles backing them up and with Ecuador still federated with the Gran Colombia, had the desire to annex Ecuador to the Republic of Peru at the expense of the Gran Colombia, feeling that Ecuador was once part of the Inca Empire.
+
+On July 28, 1821, Peruvian independence was proclaimed in Lima by the Liberator San Martin, and Tumbes and Jaen, which were included as part of the revolution of Trujillo by the Peruvian occupying force, had the whole region swear allegiance to the new Peruvian flag and incorporated itself into Peru, even though Peru was not completely liberated from Spain. After Peru was completely liberated from Spain by the patriot armies led by Bolivar and Antonio Jose de Sucre at the Battle of Ayacucho dated December 9, 1824, there was a strong desire by some Peruvians to resurrect the Inca Empire and to include Bolivia and Ecuador. One of these Peruvian Generals was the Ecuadorian-born José de La Mar, who became one of Peru's presidents after Bolivar resigned as dictator of Peru and returned to Colombia. Gran Colombia had always protested Peru for the return of Jaen and Tumbes for almost a decade, then finally Bolivar after long and futile discussion over the return of Jaen, Tumbes, and part of Mainas, declared war. President and General José de La Mar, who was born in Ecuador, believing his opportunity had come to annex the District of Ecuador to Peru, personally, with a Peruvian force, invaded and occupied Guayaquil and a few cities in the Loja region of southern Ecuador on November 28, 1828.
+
+The war ended when a triumphant heavily outnumbered southern Gran Colombian army at Battle of Tarqui dated February 27, 1829, led by Antonio José de Sucre, defeated the Peruvian invasion force led by President La Mar. This defeat led to the signing of the Treaty of Guayaquil dated September 22, 1829, whereby Peru and its Congress recognized Gran Colombian rights over Tumbes, Jaen, and Maynas. Through protocolized meetings between representatives of Peru and Gran Colombia, the border was set as Tumbes river in the west and in the east the Maranon and Amazon rivers were to be followed toward Brazil as the most natural borders between them. However, what was pending was whether the new border around the Jaen region should follow the Chinchipe River or the Huancabamba River. According to the peace negotiations Peru agreed to return Guayaquil, Tumbez, and Jaén; despite this, Peru returned Guayaquil, but failed to return Tumbes and Jaén, alleging that it was not obligated to follow the agreements, since the Gran Colombia ceased to exist when it divided itself into three different nations - Ecuador, Colombia, and Venezuela.
+
+The Central District of the Gran Colombia, known as Cundinamarca or New Granada (modern Colombia) with its capital in Bogota, did not recognize the separation of the Southern District of the Gran Colombia, with its capital in Quito, from the Gran Colombian federation on May 13, 1830. After Ecuador's separation, the Department of Cauca voluntarily decided to unite itself with Ecuador due to instability in the central government of Bogota. The Venezuelan born President of Ecuador, the general Juan José Flores, with the approval of the Ecuadorian congress annexed the Department of Cauca on December 20, 1830, since the government of Cauca had called for union with the District of the South as far back as April 1830. Moreover, the Cauca region, throughout its long history, had very strong economic and cultural ties with the people of Ecuador. Also, the Cauca region, which included such cities as Pasto, Popayán, and Buenaventura, had always been dependent on the Presidencia or Audiencia of Quito.
+
+Fruitless negotiations continued between the governments of Bogotá and Quito, where the government of Bogotá did not recognize the separation of Ecuador or that of Cauca from the Gran Colombia until war broke out in May 1832. In five months, New Granada defeated Ecuador due to the fact that the majority of the Ecuadorian Armed Forces were composed of rebellious angry unpaid veterans from Venezuela and Colombia that did not want to fight against their fellow countrymen. Seeing that his officers were rebelling, mutinying, and changing sides, President Flores had no option but to reluctantly make peace with New Granada. The Treaty of Pasto of 1832 was signed by which the Department of Cauca was turned over to New Granada (modern Colombia), the government of Bogotá recognized Ecuador as an independent country and the border was to follow the Ley de División Territorial de la República de Colombia (Law of the Division of Territory of the Gran Colombia) passed on June 25, 1824. This law set the border at the river Carchi and the eastern border that stretched to Brazil at the Caquetá river. Later, Ecuador contended that the Republic of Colombia, while reorganizing its government, unlawfully made its eastern border provisional and that Colombia extended its claims south to the Napo River because it said that the Government of Popayán extended its control all the way to the Napo River.
+
+When Ecuador seceded from the Gran Colombia, Peru decided not to follow the treaty of Guayaquil of 1829 or the protocoled agreements made. Peru contested Ecuador's claims with the newly discovered Real Cedula of 1802, by which Peru claims the King of Spain had transferred these lands from the Viceroyalty of New Granada to the Viceroyalty of Peru. During colonial times this was to halt the ever-expanding Portuguese settlements into Spanish domains, which were left vacant and in disorder after the expulsion of Jesuit missionaries from their bases along the Amazon Basin. Ecuador countered by labeling the Cedula of 1802 an ecclesiastical instrument, which had nothing to do with political borders. Peru began its de facto occupation of disputed Amazonian territories, after it signed a secret 1851 peace treaty in favor of Brazil. This treaty disregarded Spanish rights that were confirmed during colonial times by a Spanish-Portuguese treaty over the Amazon regarding territories held by illegal Portuguese settlers.
+
+Peru began occupying the defenseless missionary villages in the Mainas or Maynas region, which it began calling Loreto, with its capital in Iquitos. During its negotiations with Brazil, Peru stated that based on the royal cedula of 1802, it claimed Amazonian Basin territories up to Caqueta River in the north and toward the Andes Mountain range, depriving Ecuador and Colombia of all their claims to the Amazon Basin. Colombia protested stating that its claims extended south toward the Napo and Amazon Rivers. Ecuador protested that it claimed the Amazon Basin between the Caqueta river and the Marañon-Amazon river. Peru ignored these protests and created the Department of Loreto in 1853 with its capital in Iquitos which it had recently invaded and systematically began to occupy using the river systems in all the territories claimed by both Colombia and Ecuador. Peru briefly occupied Guayaquil again in 1860, since Peru thought that Ecuador was selling some of the disputed land for development to British bond holders, but returned Guayaquil after a few months. The border dispute was then submitted to Spain for arbitration from 1880 to 1910, but to no avail.
+
+In the early part of the 20th century, Ecuador made an effort to peacefully define its eastern Amazonian borders with its neighbours through negotiation. On May 6, 1904, Ecuador signed the Tobar-Rio Branco Treaty recognizing Brazil's claims to the Amazon in recognition of Ecuador's claim to be an Amazonian country to counter Peru's earlier Treaty with Brazil back on October 23, 1851. Then after a few meetings with the Colombian government's representatives an agreement was reached and the Muñoz Vernaza-Suarez Treaty was signed July 15, 1916, in which Colombian rights to the Putumayo river were recognized as well as Ecuador's rights to the Napo river and the new border was a line that ran midpoint between those two rivers. In this way, Ecuador gave up the claims it had to the Amazonian territories between the Caquetá River and Napo River to Colombia, thus cutting itself off from Brazil. Later, a brief war erupted between Colombia and Peru, over Peru's claims to the Caquetá region, which ended with Peru reluctantly signing the Salomon-Lozano Treaty on March 24, 1922. Ecuador protested this secret treaty, since Colombia gave away Ecuadorian claimed land to Peru that Ecuador had given to Colombia in 1916.
+
+On July 21, 1924, the Ponce-Castro Oyanguren Protocol was signed between Ecuador and Peru where both agreed to hold direct negotiations and to resolve the dispute in an equitable manner and to submit the differing points of the dispute to the United States for arbitration. Negotiations between the Ecuadorian and Peruvian representatives began in Washington on September 30, 1935. These negotiations were long and tiresome. Both sides logically presented their cases, but no one seemed to give up their claims. Then on February 6, 1937, Ecuador presented a transactional line which Peru rejected the next day. The negotiations turned into intense arguments during the next 7 months and finally on September 29, 1937, the Peruvian representatives decided to break off the negotiations without submitting the dispute to arbitration because the direct negotiations were going nowhere.
+
+Four years later in 1941, amid fast-growing tensions within disputed territories around the Zarumilla River, war broke out with Peru. Peru claimed that Ecuador's military presence in Peruvian-claimed territory was an invasion; Ecuador, for its part, claimed that Peru had recently invaded Ecuador around the Zarumilla River and that Peru since Ecuador's independence from Spain has systematically occupied Tumbez, Jaen, and most of the disputed territories in the Amazonian Basin between the Putomayo and Marañon Rivers. In July 1941, troops were mobilized in both countries. Peru had an army of 11,681 troops who faced a poorly supplied and inadequately armed Ecuadorian force of 2,300, of which only 1,300 were deployed in the southern provinces. Hostilities erupted on July 5, 1941, when Peruvian forces crossed the Zarumilla river at several locations, testing the strength and resolve of the Ecuadorian border troops. Finally, on July 23, 1941, the Peruvians launched a major invasion, crossing the Zarumilla river in force and advancing into the Ecuadorian province of El Oro.
+
+During the course of the Ecuadorian–Peruvian War, Peru gained control over part of the disputed territory and some parts of the province of El Oro, and some parts of the province of Loja, demanding that the Ecuadorian government give up its territorial claims. The Peruvian Navy blocked the port of Guayaquil, almost cutting all supplies to the Ecuadorian troops. After a few weeks of war and under pressure by the United States and several Latin American nations, all fighting came to a stop. Ecuador and Peru came to an accord formalized in the Rio Protocol, signed on January 29, 1942, in favor of hemispheric unity against the Axis Powers in World War II favouring Peru with the territory they occupied at the time the war came to an end.
+
+The 1944 Glorious May Revolution followed a military-civilian rebellion and a subsequent civic strike which successfully removed Carlos Arroyo del Río as a dictator from Ecuador's government. However, a post-Second World War recession and popular unrest led to a return to populist politics and domestic military interventions in the 1960s, while foreign companies developed oil resources in the Ecuadorian Amazon. In 1972, construction of the Andean pipeline was completed. The pipeline brought oil from the east side of the Andes to the coast, making Ecuador South America's second largest oil exporter. The pipeline in southern Ecuador did nothing to resolve tensions between Ecuador and Peru, however.
+
+The Rio Protocol failed to precisely resolve the border along a little river in the remote Cordillera del Cóndor region in southern Ecuador. This caused a long-simmering dispute between Ecuador and Peru, which ultimately led to fighting between the two countries; first a border skirmish in January–February 1981 known as the Paquisha Incident, and ultimately full-scale warfare in January 1995 where the Ecuadorian military shot down Peruvian aircraft and helicopters and Peruvian infantry marched into southern Ecuador. Each country blamed the other for the onset of hostilities, known as the Cenepa War. Sixto Durán Ballén, the Ecuadorian president, famously declared that he would not give up a single centimeter of Ecuador. Popular sentiment in Ecuador became strongly nationalistic against Peru: graffiti could be seen on the walls of Quito referring to Peru as the "Cain de Latinoamérica", a reference to the murder of Abel by his brother Cain in the Book of Genesis.[27]
+
+Ecuador and Peru signed the Brasilia Presidential Act peace agreement on October 26, 1998, which ended hostilities, and effectively put an end to the Western Hemisphere's longest running territorial dispute.[28] The Guarantors of the Rio Protocol (Argentina, Brazil, Chile, and the United States of America) ruled that the border of the undelineated zone was to be set at the line of the Cordillera del Cóndor. While Ecuador had to give up its decades-old territorial claims to the eastern slopes of the Cordillera, as well as to the entire western area of Cenepa headwaters, Peru was compelled to give to Ecuador, in perpetual lease but without sovereignty, 1 km2 (0.39 sq mi) of its territory, in the area where the Ecuadorian base of Tiwinza – focal point of the war – had been located within Peruvian soil and which the Ecuadorian Army held during the conflict. The final border demarcation came into effect on May 13, 1999 and the multi-national MOMEP (Military Observer Mission for Ecuador and Peru) troop deployment withdrew on June 17, 1999.[28]
+
+In 1972, a "revolutionary and nationalist" military junta overthrew the government of Velasco Ibarra. The coup d'état was led by General Guillermo Rodríguez and executed by navy commander Jorge Queirolo G. The new president exiled José María Velasco to Argentina. He remained in power until 1976, when he was removed by another military government. That military junta was led by Admiral Alfredo Poveda, who was declared chairman of the Supreme Council. The Supreme Council included two other members: General Guillermo Durán Arcentales and General Luis Leoro Franco. The civil society more and more insistently called for democratic elections. Colonel Richelieu Levoyer, Government Minister, proposed and implemented a Plan to return to the constitutional system through universal elections. This plan enabled the new democratically elected president to assume the duties of the executive office.
+
+Elections were held on April 29, 1979, under a new constitution. Jaime Roldós Aguilera was elected president, garnering over one million votes, the most in Ecuadorian history. He took office on August 10, as the first constitutionally elected president after nearly a decade of civilian and military dictatorships. In 1980, he founded the Partido Pueblo, Cambio y Democracia (People, Change, and Democracy Party) after withdrawing from the Concentración de Fuerzas Populares (Popular Forces Concentration) and governed until May 24, 1981, when he died along with his wife and the minister of defense, Marco Subia Martinez, when his Air Force plane crashed in heavy rain near the Peruvian border. Many people believe that he was assassinated by the CIA,[citation needed] given the multiple death threats leveled against him because of his reformist agenda, deaths in automobile crashes of two key witnesses before they could testify during the investigation, and the sometimes contradictory accounts of the incident.
+
+Roldos was immediately succeeded by Vice President Osvaldo Hurtado, who was followed in 1984 by León Febres Cordero from the Social Christian Party. Rodrigo Borja Cevallos of the Democratic Left (Izquierda Democrática, or ID) party won the presidency in 1988, running in the runoff election against Abdalá Bucaram (brother in law of Jaime Roldos and founder of the Ecuadorian Roldosist Party). His government was committed to improving human rights protection and carried out some reforms, notably an opening of Ecuador to foreign trade. The Borja government concluded an accord leading to the disbanding of the small terrorist group, "¡Alfaro Vive, Carajo!" ("Alfaro Lives, Dammit!"), named after Eloy Alfaro. However, continuing economic problems undermined the popularity of the ID, and opposition parties gained control of Congress in 1999.
+
+The emergence of the Amerindian population as an active constituency has added to the democratic volatility of the country in recent years. The population has been motivated by government failures to deliver on promises of land reform, lower unemployment and provision of social services, and historical exploitation by the land-holding elite. Their movement, along with the continuing destabilizing efforts by both the elite and leftist movements, has led to a deterioration of the executive office. The populace and the other branches of government give the president very little political capital, as illustrated by the most recent removal of President Lucio Gutiérrez from office by Congress in April 2005. Vice President Alfredo Palacio took his place and remained in office until the presidential election of 2006, in which Rafael Correa gained the presidency.[29]
+
+In December 2008, president Correa declared Ecuador's national debt illegitimate, based on the argument that it was odious debt contracted by corrupt and despotic prior regimes. He announced that the country would default on over $3 billion worth of bonds; he then pledged to fight creditors in international courts and succeeded in reducing the price of outstanding bonds by more than 60%.[30] He brought Ecuador into the Bolivarian Alliance for the Americas in June 2009. To date, Correa's administration has succeeded in reducing the high levels of poverty and unemployment in Ecuador.[31][32][33][34][35]
+
+After being elected in 2017, President Lenin Moreno's government adopted economically liberal policies: reduction of public spending, trade liberalization, flexibility of the labour code, etc. He also left the left-wing Bolivarian Alliance for the Americas in August 2018. The Productive Development Act enshrines an austerity policy, and reduces the development and redistribution policies of the previous mandate. In the area of taxes, the authorities aim to "encourage the return of investors" by granting amnesty to fraudsters and proposing measures to reduce tax rates for large companies. In addition, the government waives the right to tax increases in raw material prices and foreign exchange repatriations.[36]
+
+A  series of protests began on 3 October 2019 against the end of fuel subsidies and austerity measures adopted by President of Ecuador Lenín Moreno and his administration. On 10 October, protesters overran the capital Quito causing the Government of Ecuador to relocate to Guayaquil,[37] but it was reported that the government still had plans to return to Quito.[38]
+
+The Ecuadorian State consists of five branches of government: the Executive Branch, the Legislative Branch, the Judicial Branch, the Electoral Branch, and Transparency and Social Control.
+
+Ecuador is governed by a democratically elected president, for a four-year term. The current president of Ecuador, Lenín Moreno, exercises his power from the presidential Palacio de Carondelet in Quito. The current constitution was written by the Ecuadorian Constituent Assembly elected in 2007, and was approved by referendum in 2008. Since 1936, voting is compulsory for all literate persons aged 18–65, optional for all other citizens.[39]
+
+The executive branch includes 23 ministries. Provincial governors and councilors (mayors, aldermen, and parish boards) are directly elected. The National Assembly of Ecuador meets throughout the year except for recesses in July and December. There are thirteen permanent committees. Members of the National Court of Justice are appointed by the National Judicial Council for nine-year terms.
+
+The executive branch is led by the president, an office currently held by Lenín Moreno. He is accompanied by the vice-president, currently Otto Sonnenholzner, elected for four years (with the ability to be re-elected only once). As head of state and chief government official, he is responsible for public administration including the appointing of national coordinators, ministers, ministers of State and public servants. The executive branch defines foreign policy, appoints the Chancellor of the Republic, as well as ambassadors and consuls, being the ultimate authority over the Armed Forces of Ecuador, National Police of Ecuador, and appointing authorities. The acting president's wife receives the title of First Lady of Ecuador.
+
+The legislative branch is embodied by the National Assembly, which is headquartered in the city of Quito in the Legislative Palace, and consists of 137 assemblymen, divided into ten committees and elected for a four-year term. Fifteen national constituency elected assembly, two Assembly members elected from each province and one for every 100,000 inhabitants or fraction exceeding 150,000, according to the latest national population census. In addition, statute determines the election of assembly of regions and metropolitan districts.
+
+Ecuador's judiciary has as its main body the Judicial Council, and also includes the National Court of Justice, provincial courts, and lower courts. Legal representation is made by the Judicial Council.
+The National Court of Justice is composed of 21 judges elected for a term of nine years. Judges are renewed by thirds every three years pursuant to the Judicial Code. These are elected by the Judicial Council on the basis of opposition proceedings and merits.
+The justice system is buttressed by the independent offices of public prosecutor and the public defender. Auxiliary organs are as follows: notaries, court auctioneers, and court receivers. Also there is a special legal regime for Amerindians.
+
+The electoral system functions by authorities which enter only every four years or when elections or referendums occur. Its main functions are to organize, control elections, and punish the infringement of electoral rules. Its main body is the National Electoral Council, which is based in the city of Quito, and consists of seven members of the political parties most voted, enjoying complete financial and administrative autonomy. This body, along with the electoral court, forms the Electoral Branch which is one of Ecuador's five branches of government.
+
+The Transparency and Social Control consists of the Council of Citizen Participation and Social Control, an ombudsman, the Comptroller General of the State, and the superintendents. Branch members hold office for five years. This branch is responsible for promoting transparency and control plans publicly, as well as plans to design mechanisms to combat corruption, as also designate certain authorities, and be the regulatory mechanism of accountability in the country.
+
+UN's Human Rights Council's (HRC) Universal Periodic Review (UPR) has treated the restrictions on freedom of expression and efforts to control NGOs and recommended that Ecuador should stop the criminal sanctions for the expression of opinions, and delay in implementing judicial reforms. Ecuador rejected the recommendation on decriminalization of libel.[40]
+
+According to Human Rights Watch (HRW) President Correa has intimidated journalists and subjected them to "public denunciation and retaliatory litigation". The sentences to journalists have been years of imprisonment and millions of dollars of compensation, even though defendants have been pardoned.[40] Correa has stated he was only seeking a retraction for slanderous statements.[41]
+
+According to HRW, Correa's government has weakened the freedom of press and independence of the judicial system. In Ecuador's current judicial system, judges are selected in a contest of merits, rather than government appointments. However, the process of selection has been criticized as biased and subjective. In particular, the final interview is said to be given "excessive weighing". Judges and prosecutors that have made decisions in favor of Correa in his lawsuits have received permanent posts, while others with better assessment grades have been rejected.[40][42]
+
+The laws also forbid articles and media messages that could favor or disfavor some political message or candidate. In the first half of 2012, twenty private TV or radio stations were closed down.[40]
+
+In July 2012, the officials warned the judges that they would be sanctioned and possibly dismissed if they allowed the citizens to appeal to the protection of their constitutional rights against the state.[40]
+
+People engaging in public protests against environmental and other issues are prosecuted for "terrorism and sabotage", which may lead to an eight-year prison sentence.[40]
+
+Ecuador joined the Organization of Petroleum Exporting Countries (OPEC) in 1973 and suspended its membership in 1992. Under President Rafael Correa, the country returned to OPEC before leaving again in 2020 under the instruction of President Moreno, citing its desire to increase crude oil importation to gain more revenue.[43][44]
+
+In Antarctica, Ecuador has maintained a peaceful research station for scientific study as a member nation of the Antarctica Treaty. Ecuador has often placed great emphasis on multilateral approaches to international issues. Ecuador is a member of the United Nations (and most of its specialized agencies) and a member of many regional groups, including the Rio Group, the Latin American Economic System, the Latin American Energy Organization, the Latin American Integration Association, the Andean Community of Nations, and the Bank of the South (Spanish: Banco del Sur or BancoSur).
+
+In 2017, the Ecuadorian parliament adopted a Law on human mobility.[45]
+
+The International Organization for Migration lauds Ecuador as the first state to have established the promotion of the concept of universal citizenship in its constitution, aiming to promote the universal recognition and protection of the human rights of migrants.[46] In 2017, Ecuador signed the UN treaty on the Prohibition of Nuclear Weapons.[47]
+
+Ecuador is divided into 24 provinces (Spanish: provincias), each with its own administrative capital:
+
+The provinces are divided into cantons and further subdivided into parishes (parroquias).
+
+Regionalization, or zoning, is the union of two or more adjoining provinces in order to decentralize the administrative functions of the capital, Quito.
+In Ecuador, there are seven regions, or zones, each shaped by the following provinces:
+
+Quito and Guayaquil are Metropolitan Districts. Galápagos, despite being included within Region 5,[49] is also under a special unit.[50]
+
+The Ecuadorian Armed Forces (Fuerzas Armadas de la Republica de Ecuador), consists of the Army, Air Force, and Navy and have the stated responsibility for the preservation of the integrity and national sovereignty of the national territory.
+
+The military tradition starts in Gran Colombia, where a sizable army was stationed in Ecuador due to border disputes with Peru, which claimed territories under its political control when it was a Spanish vice-royalty. Once Gran Colombia was dissolved after the death of Simón Bolívar in 1830, Ecuador inherited the same border disputes and had the need of creating its own professional military force. So influential was the military in Ecuador in the early republican period that its first decade was under the control of General Juan José Flores, first president of Ecuador of Venezuelan origin. General Jose Ma. Urbina and General Robles are examples of military figures who became presidents of the country in the early republican period.
+
+Due to the continuous border disputes with Peru, finally settled in the early 2000s, and due to the ongoing problem with the Colombian guerrilla insurgency infiltrating Amazonian provinces, the Ecuadorian Armed Forces has gone through a series of changes. In 2009, the new administration at the Defense Ministry launched a deep restructuring within the forces, increasing spending budget to $1,691,776,803, an increase of 25%.[51]
+
+The icons of the Ecuadorian military forces are Marshall Antonio José de Sucre and General Eloy Alfaro.
+
+The Military Academy General Eloy Alfaro (c. 1838) located in Quito is in charge to graduate the army officers.[52]
+
+The Ecuadorian Navy Academy (c. 1837), located in Salinas graduates the navy officers.[53]
+
+The Air Academy "Cosme Rennella (c. 1920), also located in Salinas, graduates the air force officers.[54]
+
+Other training academies for different military specialties are found across the country.
+
+The IWIAS is a special force trained to perform exploration and militar activities. Is considered the best elite force of Ecuador and is conformed by indigenous of the Amazon who combine their inherital experience for jungle dominance with modern army tactics.
+
+Ecuador has a total area of 283,561 km2 (109,484 sq mi), including the Galápagos Islands. Of this, 276,841 km2 (106,889 sq mi) is land and 6,720 km2 (2,595 sq mi) water.[1] Ecuador is bigger than Uruguay, Suriname, Guyana and French Guyana in South America.
+
+Ecuador lies between latitudes 2°N and 5°S,
+bounded on the west by the Pacific Ocean, and has 2,337 km (1,452 mi) of coastline. It has 2,010 km (1,250 mi) of land boundaries, with Colombia in the north (with a 590 km (367 mi) border) and Peru in the east and south (with a 1,420 km (882 mi) border). It is the westernmost country that lies on the equator.[55]
+
+The country has four main geographic regions:
+
+Ecuador's capital is Quito, which is in the province of Pichincha in the Sierra region. Its largest city is Guayaquil, in the Guayas Province. Cotopaxi, just south of Quito, is one of the world's highest active volcanoes. The top of Mount Chimborazo (6,268 m, or 20,560 ft, above sea level), Ecuador's tallest mountain, is the most distant point from the center of the Earth on the Earth's surface because of the ellipsoid shape of the planet.[1]
+
+There is great variety in the climate, largely determined by altitude. It is mild year-round in the mountain valleys, with a humid subtropical climate in coastal areas and rainforest in lowlands. The Pacific coastal area has a tropical climate with a severe rainy season. The climate in the Andean highlands is temperate and relatively dry, and the Amazon basin on the eastern side of the mountains shares the climate of other rainforest zones.
+
+Because of its location at the equator, Ecuador experiences little variation in daylight hours during the course of a year. Both sunrise and sunset occur each day at the two six o'clock hours.[1]
+
+The Andes is the watershed divisor between the Amazon watershed, which runs to the east, and the Pacific, including the north–south rivers Mataje, Santiago, Esmeraldas, Chone, Guayas, Jubones, and Puyango-Tumbes.
+
+Almost all of the rivers in Ecuador form in the Sierra region and flow east toward the Amazon River or west toward the Pacific Ocean. The rivers rise from snowmelt at the edges of the snowcapped peaks or from the abundant precipitation that falls at higher elevations. In the Sierra region, the streams and rivers are narrow and flow rapidly over precipitous slopes. Rivers may slow and widen as they cross the hoyas yet become rapid again as they flow from the heights of the Andes to the lower elevations of the other regions. The highland rivers broaden as they enter the more level areas of the Costa and the Oriente.
+
+In the Costa, the external coast has mostly intermittent rivers that are fed by constant rains from December through May and become empty riverbeds during the dry season. The few exceptions are the longer, perennial rivers that flow throughout the external coast from the internal coast and La Sierra on their way to the Pacific Ocean. The internal coast, by contrast, is crossed by perennial rivers that may flood during the rainy season, sometimes forming swamps.
+
+Major rivers in the Oriente include the Pastaza, Napo, and Putumayo. The Pastaza is formed by the confluence of the Chambo and the Patate rivers, both of which rise in the Sierra. The Pastaza includes the Agoyan waterfall, which at sixty-one meters (200 feet) is the highest waterfall in Ecuador. The Napo rises near Mount Cotopaxi and is the major river used for transport in the eastern lowlands. The Napo ranges in width from 500 to 1,800 m (1,640 to 5,906 ft). In its upper reaches, the Napo flows rapidly until the confluence with one of its major tributaries, the Coca River, where it slows and levels off. The Putumayo forms part of the border with Colombia. All of these rivers flow into the Amazon River. The Galápagos Islands have no significant rivers. Several of the larger islands, however, have freshwater springs although they are surrounded by the Pacific Ocean.
+
+Ecuador is one of seventeen megadiverse countries in the world according to Conservation International,[19] and it has the most biodiversity per square kilometer of any nation.[56][57]
+
+Ecuador has 1,600 bird species (15% of the world's known bird species) in the continental area and 38 more endemic in the Galápagos. In addition to over 16,000 species of plants, the country has 106 endemic reptiles, 138 endemic amphibians, and 6,000 species of butterfly. The Galápagos Islands are well known as a region of distinct fauna, famous as the place of birth of Darwin's Theory of Evolution and a UNESCO World Heritage Site.[58]
+
+Ecuador has the first constitution to recognize the rights of nature.[59] The protection of the nation's biodiversity is an explicit national priority as stated in the National Plan of "Buen Vivir", or good living, Objective 4, "Guarantee the rights of nature", Policy 1: "Sustainably conserve and manage the natural heritage, including its land and marine biodiversity, which is considered a strategic sector".[56]  As of the writing of the Plan in 2008, 19% of Ecuador's land area was in a protected area; however, the Plan also states that 32% of the land must be protected in order to truly preserve the nation's biodiversity.[56] Current protected areas include 11 national parks, 10 wildlife refuges, 9 ecological reserves, and other areas.[60] A program begun in 2008, Sociobosque, is preserving another 2.3% of total land area (6,295 km2, or 629,500 ha) by paying private landowners or community landowners (such as Amerindian tribes) incentives to maintain their land as native ecosystems such as native forests or grasslands. Eligibility and subsidy rates for this program are determined based on the poverty in the region, the number of hectares that will be protected, and the type of ecosystem of the land to be protected, among other factors.[61]
+
+Despite being on the UNESCO list, the Galápagos are endangered by a range of negative environmental effects, threatening the existence of this exotic ecosystem.[62] Additionally, oil exploitation of the Amazon rainforest has led to the release of billions of gallons of untreated wastes, gas, and crude oil into the environment, contaminating ecosystems and causing detrimental health effects to Amerindian peoples.[63]
+
+Ecuador has a developing economy that is highly dependent on commodities, namely petroleum and agricultural products. The country is classified as an upper-middle-income country. Ecuador's economy is the eighth largest in Latin America and experienced an average growth of 4.6% between 2000 and 2006.[64][failed verification] From 2007 to 2012, Ecuador's GDP grew at an annual average of 4.3 percent, above the average for Latin America and the Caribbean, which was 3.5%, according to the United Nations' Economic Commission for Latin American and the Caribbean (ECLAC).[65] Ecuador was able to maintain relatively superior growth during the crisis. In January 2009, the Central Bank of Ecuador (BCE) put the 2010 growth forecast at 6.88%.[66] In 2011, its GDP grew at 8% and ranked 3rd highest in Latin America, behind Argentina (2nd) and Panama (1st).[67] Between 1999 and 2007, GDP doubled, reaching $65,490 million according to BCE.[68]
+The inflation rate until January 2008, was about 1.14%, the highest in the past year, according to the government.[69][70] The monthly unemployment rate remained at about 6 and 8 percent from December 2007 until September 2008; however, it went up to about 9 percent in October and dropped again in November 2008 to 8 percent.[71] Unemployment mean annual rate for 2009 in Ecuador was 8.5% because the global economic crisis continued to affect the Latin American economies. From this point, unemployment rates started a downward trend: 7.6% in 2010, 6.0% in 2011, and 4.8% in 2012.[72]
+
+The extreme poverty rate has declined significantly between 1999 and 2010.[73] In 2001, it was estimated at 40% of the population, while by 2011 the figure dropped to 17.4% of the total population.[74] This is explained to an extent by emigration and the economic stability achieved after adopting the U.S. dollar as official means of transaction (before 2000, the Ecuadorian sucre was prone to rampant inflation). However, starting in 2008, with the bad economic performance of the nations where most Ecuadorian emigrants work, the reduction of poverty has been realized through social spending, mainly in education and health.[75]
+
+Oil accounts for 40% of exports and contributes to maintaining a positive trade balance.[76] Since the late 1960s, the exploitation of oil increased production, and proven reserves are estimated at 6.51 billion barrels as of 2011[update].[77]
+
+The overall trade balance for August 2012 was a surplus of almost $390 million for the first six months of 2012, a huge figure compared with that of 2007, which reached only $5.7 million; the surplus had risen by about $425 million compared to 2006.[74] The oil trade balance positive had revenues of $3.295 million in 2008, while non-oil was negative, amounting to $2.842 million. The trade balance with the United States, Chile, the European Union, Bolivia, Peru, Brazil, and Mexico is positive. The trade balance with Argentina, Colombia, and Asia is negative.[78]
+
+In the agricultural sector, Ecuador is a major exporter of bananas (first place worldwide in production and export), flowers, and the seventh largest producer of cocoa.[79] Ecuador also produces coffee, rice, potatoes, cassava (manioc, tapioca), plantains and sugarcane; cattle, sheep, pigs, beef, pork and dairy products; fish, and shrimp; and balsa wood.[80] The country's vast resources include large amounts of timber across the country, like eucalyptus and mangroves.[81] Pines and cedars are planted in the region of La Sierra and walnuts, rosemary, and balsa wood in the Guayas River Basin.[82]
+The industry is concentrated mainly in Guayaquil, the largest industrial center, and in Quito, where in recent years the industry has grown considerably. This city is also the largest business center of the country.[83] Industrial production is directed primarily to the domestic market.[citation needed] Despite this, there is limited export of products produced or processed industrially.[citation needed] These include canned foods, liquor, jewelry, furniture, and more.[citation needed] A minor industrial activity is also concentrated in Cuenca.[84] Incomes from tourism has been increasing during the last few years because of the Government showing the variety of climates and the biodiversity of Ecuador.
+
+Ecuador has negotiated bilateral treaties with other countries, besides belonging to the Andean Community of Nations,[85] and an associate member of Mercosur.[86] It also serves on the World Trade Organization (WTO), in addition to the Inter-American Development Bank (IDB), World Bank, International Monetary Fund (IMF), Corporación Andina de Fomento (CAF) and other multilateral agencies.[87][88][89] In April 2007, Ecuador paid off its debt to the IMF, thus ending an era of interventionism of the Agency in the country.[90][91]
+The public finance of Ecuador consists of the Central Bank of Ecuador (BCE), the National Development Bank (BNF), the State Bank.
+
+The Ministry of Information and Tourism was created on August 10, 1992, at the beginning of the government of Sixto Durán Ballén, who viewed tourism as a fundamental activity for the economic and social development of the peoples. Faced with the growth of the tourism sector, in June 1994, the decision was taken to separate tourism from information, so that it is exclusively dedicated to promoting and strengthening this activity.
+
+Ecuador is a country with vast natural wealth. The diversity of its four regions has given rise to thousands of species of flora and fauna. It has around 1640 kinds of birds. The species of butterflies border the 4,500, the reptiles 345, the amphibians 358 and the mammals 258, among others. Not in vain, Ecuador is considered one of the 17 countries where the planet's highest biodiversity is concentrated, being also the largest country with diversity per km2 in the world. Most of its fauna and flora lives in 26 protected areas by the State. Also, it has a huge culture spectrum. Since 2007, with the government of Rafael Correa, the tourism brand "Ecuador Ama la Vida" has been transformed, with which the nation's tourism promotion would be sold. Focused on considering it as a country friendly and respectful of the nature, natural biodiversity and cultural diversity of the peoples. And for this, means of exploiting them are developed along with the private economy.
+
+The country has two cities UNESCO World Heritage Sites: Quito and Cuenca, as well as two natural UNESCO World Heritage Sites: the Galapagos Islands and Sangay National Park in addition to one World Biosphere Reserve, such as the Cajas Massif. Culturally, the Toquilla straw hat and the culture of the Zapara indigenous people are recognized. The most popular sites for national and foreign tourists have different nuances due to the various tourist activities offered by the country.
+
+Among the main tourist destinations are:
+
+The rehabilitation and reopening of the Ecuadorian railroad and use of it as a tourist attraction is one of the recent developments in transportation matters.[93]
+
+The roads of Ecuador in recent years have undergone important improvement. The major routes are Pan American (under enhancement from four to six lanes from Rumichaca to Ambato, the conclusion of 4 lanes on the entire stretch of Ambato and Riobamba and running via Riobamba to Loja). In the absence of the section between Loja and the border with Peru, there are the Route Espondilus and/or Ruta del Sol (oriented to travel along the Ecuadorian coastline) and the Amazon backbone (which crosses from north to south along the Ecuadorian Amazon, linking most and more major cities of it).
+
+Another major project is developing the road Manta – Tena, the highway Guayaquil – Salinas Highway Aloag Santo Domingo, Riobamba – Macas (which crosses Sangay National Park). Other new developments include the National Unity bridge complex in Guayaquil, the bridge over the Napo river in Francisco de Orellana, the Esmeraldas River Bridge in the city of the same name, and, perhaps the most remarkable of all, the Bahia – San Vincente Bridge, being the largest on the Latin American Pacific coast.
+
+Cuenca's tramway is the largest public transport system in the city and the first modern tramway in Ecuador. It was inaugurated on March 8, 2019. It has 20,4 km and 27 stations. It will transport 120 000 passagers daily. Its route starts in the south of Cuenca and ends in the north at the Parque Industrial neighbourhood.
+
+The Mariscal Sucre International Airport in Quito and the José Joaquín de Olmedo International Airport in Guayaquil have experienced a high increase in demand and have required modernization. In the case of Guayaquil it involved a new air terminal, once considered the best in South America and the best in Latin America[94] and in Quito where an entire new airport has been built in Tababela and was inaugurated in February 2013, with Canadian assistance. However, the main road leading from Quito city centre to the new airport will only be finished in late 2014, making current travelling from the airport to downtown Quito as long as two hours during rush hour.[95] Quito's old city-centre airport is being turned into parkland, with some light industrial use.
+
+Ecuador's population is ethnically diverse and the 2018 estimates put Ecuador's population at 17,084,358.[6][7] The largest ethnic group (as of 2010[update]) is the Mestizos, who are the descendants of Spanish colonists that interbred with Amerindian peoples, and constitute about 71% of the population. The White Ecuadorians (White Latin American) are a minority accounting for 6.1% of the population of Ecuador and can be found throughout all of Ecuador primarily around the urban areas. Even though Ecuador's white population during its colonial era were mainly descendants from Spain, today Ecuador's white population is a result of a mixture of European immigrants, predominantly from Spain with people from Italy, Germany, France, and Switzerland who have settled in the early 20th century. Ecuador also has people of middle eastern extraction that have also joined the ranks of the white minority. These include economically well off immigrants of Lebanese and Palestinian descent, who are either Christian or Muslim (Islam in Ecuador). In addition, there is a small European Jewish (Ecuadorian Jews) population, which is based mainly in Quito and to a lesser extent in Guayaquil.[3] Amerindians account for 7% of the current population. The mostly rural Montubio population of the coastal provinces of Ecuador, who might be classified as Pardo account for 7.4% of the population. The Afro-Ecuadorians are a minority population (7%) in Ecuador, that includes the Mulattos and zambos, and are largely based in the Esmeraldas province and to a lesser degree in the predominantly Mestizo provinces of Coastal Ecuador - Guayas and Manabi. In the Highland Andes where a predominantly Mestizo, white and Amerindian population exist, the African presence is almost non-existent except for a small community in the province of Imbabura called Chota Valley.
+
+According to the Ecuadorian National Institute of Statistics and Census, 91.95% of the country's population have a religion, 7.94% are atheists and 0.11% are agnostics. Among the people that have a religion, 80.44% are Roman Catholic Latin Rite (see List of Roman Catholic dioceses in Ecuador), 11.30% are Evangelical Protestants, 1.29% are Jehovah's Witnesses and 6.97% other (mainly Jewish, Buddhists and Latter-day Saints).[97][98]
+
+In the rural parts of Ecuador, Amerindian beliefs and Catholicism are sometimes syncretized. Most festivals and annual parades are based on religious celebrations, many incorporating a mixture of rites and icons.[citation needed]
+
+There is a small number of Eastern Orthodox Christians, Amerindian religions, Muslims (see Islam in Ecuador), Buddhists and Bahá'í. According to their own estimates, The Church of Jesus Christ of Latter-day Saints accounts for about 1.4% of the population, or 211,165 members at the end of 2012.[99] According to their own sources, in 2017 there were 92,752 Jehovah's Witnesses in the country.[100]
+
+The first Jews arrived in Ecuador in the 16th and 17th centuries. Most of them are Sephardic Anusim (Crypto-Jews) and many still speak Judaeo-Spanish (Ladino) language.[101][citation needed] Today the Jewish Community of Ecuador (Comunidad Judía del Ecuador) has its seat in Quito and has approximately 200 members. Nevertheless, this number is declining because young people leave the country for the United States or Israel. The Community has a Jewish Center with a synagogue, a country club, and a cemetery. It supports the "Albert Einstein School", where Jewish history, religion, and Hebrew classes are offered. There are very small communities in Cuenca. The "Comunidad de Culto Israelita" reunites the Jews of Guayaquil. This community works independently from the "Jewish Community of Ecuador" and is composed of only 30 people.[102]
+
+Iglesia de San Sebastián church in Cuenca
+
+Iglesia y Convento de San Francisco in Quito
+
+The Ecuadorian constitution recognizes the "pluri-nationality" of those who want to exercise their affiliation with their native ethnic groups. Thus, in addition to criollos, mestizos, and Afro-Ecuadorians, some people belong to the Amerindian nations scattered in a few places in the coast, Quechua Andean villages, and the Amazonian jungle.
+
+According to a 2015 genealogical DNA testing, the average Ecuadorian is estimated to be 52.96% Native American, 41.77% European, and 5.26% Sub-Saharan African overall.[104]
+
+The majority of Ecuadorians live in the central provinces, the Andes mountains, or along the Pacific coast. The tropical forest region to the east of the mountains (El Oriente) remains sparsely populated and contains only about 3% of the population. Birth rate is 2-1 for each death. Marriages are usually from 14 and above using parental consent. About 12.4% of the population is married in the ages 15–19. Divorce rates are moderate.
+
+The five largest cities in the country are Quito (2.78 million inhabitants), Guayaquil (2.72 million inhabitants), Cuenca (636,996 inhabitants), Santo Domingo (458,580 inhabitants), and Ambato (387,309 inhabitants). While the most populated metropolitan areas of the country are those of Guayaquil, Quito, Cuenca, Manabí Centro (Portoviejo-Manta) and Ambato.[105]
+
+A small East Asian Latino community, estimated at 2,500, mainly consists of those of Japanese and Chinese descent, whose ancestors arrived as miners, farmhands and fishermen in the late 19th century.[1]
+
+In the early years of World War II, Ecuador still admitted a certain number of immigrants, and in 1939, when several South American countries refused to accept 165 Jewish refugees from Germany aboard the ship Koenigstein, Ecuador granted them entry permits.[107]
+
+In recent years, Ecuador has grown in popularity among North American expatriates.[108] They're drawn there by the authentic cultural experience and beautiful natural surroundings. Also, Ecuador's favorable residency options make for an easy transition for those who decide to settle there indefinitely.
+
+Another perk that draws many expats to Ecuador is its low cost of living. Since everything from gas to groceries costs far less than in North America, it is a popular choice for those who are looking to make the most of their retirement budget.[109]
+
+Even real estate in Ecuador is much less than its tropical counterparts. However, as more and more North Americans are discovering Ecuador's potential, property prices are beginning to rise from where they were a decade ago, particularly in the areas that are popular among expats and tourists.
+
+Ecuador's mainstream culture is defined by its Hispanic mestizo majority, and, like their ancestry, it is traditionally of Spanish heritage, influenced in different degrees by Amerindian traditions and in some cases by African elements. The first and most substantial wave of modern immigration to Ecuador consisted of Spanish colonists, following the arrival of Europeans in 1499. A lower number of other Europeans and North Americans migrated to the country in the late 19th and early twentieth centuries and, in smaller numbers, Poles, Lithuanians, English, Irish, and Croats during and after the Second World War.
+
+Since African slavery was not the workforce of the Spanish colonies in the Andes Mountains, given the subjugation of the Amerindian people through proselytization and encomiendas, the minority population of African descent is mostly found in the coastal northern province of Esmeraldas. This is largely owing to the 17th-century shipwreck of a slave-trading galleon off the northern coast of Ecuador. The few black African survivors swam to the shore and penetrated the then-thick jungle under the leadership of Anton, the chief of the group, where they remained as free men maintaining their original culture, not influenced by the typical elements found in other provinces of the coast or in the Andean region. A little later, freed slaves from Colombia known as cimarrones joined them. In the small Chota Valley of the province of Imbabura exists a small community of Africans among the province's predominantly mestizo population. These blacks are descendants of Africans, who were brought over from Colombia by Jesuits to work their colonial sugar plantations as slaves. As a general rule, small elements of zambos and mulattoes coexisted among the overwhelming mestizo population of coastal Ecuador throughout its history as gold miners in Loja, Zaruma, and Zamora and as shipbuilders and plantation workers around the city of Guayaquil. Today you can find a small community of Africans in the Catamayo valley of the predominantly mestizo population of Loja.
+
+Ecuador's Amerindian communities are integrated into the mainstream culture to varying degrees,[110] but some may also practice their own native cultures, particularly the more remote Amerindian communities of the Amazon basin. Spanish is spoken as the first language by more than 90% of the population and as a first or second language by more than 98%. Part of Ecuador's population can speak Amerindian languages, in some cases as a second language. Two percent of the population speak only Amerindian languages.
+
+Most Ecuadorians speak Spanish as their first language, with its ubiquity permeating and dominating most of the country, though there are many who speak an Amerindian language, such as Kichwa (also spelled Quechua), which is one of the Quechuan languages and is spoken by approximately 2.5 million people in Ecuador, Bolivia, Colombia, and Peru.[111] Other Amerindian languages spoken in Ecuador include Awapit (spoken by the Awá), A'ingae (spoken by the Cofan), Shuar Chicham (spoken by the Shuar), Achuar-Shiwiar (spoken by the Achuar and the Shiwiar), Cha'palaachi (spoken by the Chachi), Tsa'fiki (spoken by the Tsáchila), Paicoca (spoken by the Siona and Secoya), and Wao Tededeo (spoken by the Waorani). Use of these Amerindian languages are, however, gradually diminishing due to Spanish's widespread use in education. Though most features of Ecuadorian Spanish are universal to the Spanish-speaking world, there are several idiosyncrasies.
+
+The music of Ecuador has a long history. Pasillo is a genre of indigenous Latin music. In Ecuador it is the "national genre of music". Through the years, many cultures have brought their influences together to create new types of music. There are also different kinds of traditional music like albazo, pasacalle, fox incaico, tonada, capishca, Bomba (highly established in Afro-Ecuadorian societies), and so on. Tecnocumbia and Rockola are clear examples of the influence of foreign cultures. One of the most traditional forms of dancing in Ecuador is Sanjuanito. It is originally from northern Ecuador (Otavalo-Imbabura). Sanjuanito is a type of dance music played during festivities by the mestizo and Amerindian communities. According to the Ecuadorian musicologist Segundo Luis Moreno, Sanjuanito was danced by Amerindian people during San Juan Bautista's birthday. This important date was established by the Spaniards on June 24, coincidentally the same date when Amerindian people celebrated their rituals of Inti Raymi.
+
+Ecuadorian cuisine is diverse, varying with the altitude and associated agricultural conditions. Most regions in Ecuador follow the traditional three-course meal of soup, a course that includes rice and a protein, and then dessert and coffee to finish. Supper is usually lighter and sometimes consists only of coffee or herbal tea with bread.[citation needed]
+
+In the highland region; grilled pork, chicken, beef, and cuy (guinea pig) are popular and are served with a variety of grains (especially rice and mote) or potatoes.[citation needed]
+
+In the coastal region, seafood is very popular, with fish, shrimp, and ceviche being key parts of the diet. Generally, ceviches are served with fried plantain (chifles or patacones), popcorn, or tostado. Plantain- and peanut-based dishes are the basis of most coastal meals. Encocados (dishes that contain a coconut sauce) are also very popular. Churrasco is a staple food of the coastal region, especially Guayaquil. Arroz con menestra y carne asada (rice with beans and grilled beef) is one of the traditional dishes of Guayaquil, as is fried plantain, which is often served with it. This region is a leading producer of bananas, cocoa beans (to make chocolate), shrimp, tilapia, mango, and passion fruit, among other products.[citation needed]
+
+In the Amazon region, a dietary staple is the yuca, elsewhere called cassava. Many fruits are available in this region, including bananas, tree grapes, and peach palms.[citation needed]
+
+Early literature in colonial Ecuador, as in the rest of Spanish America, was influenced by the Spanish Golden Age. One of the earliest examples is Jacinto Collahuazo,[112] an Amerindian chief of a northern village in today's Ibarra, born in the late 1600s. Despite the early repression and discrimination of the native people by the Spanish, Collahuazo learned to read and write in Castilian, but his work was written in Quechua. The use of Quipu was banned by the Spanish,[113] and in order to preserve their work, many Inca poets had to resort to the use of the Latin alphabet to write in their native Quechua language. The history behind the Inca drama "Ollantay", the oldest literary piece in existence for any Amerindian language in America,[114] shares some similarities with the work of Collahuazo. Collahuazo was imprisoned and all of his work burned. The existence of his literary work came to light many centuries later, when a crew of masons was restoring the walls of a colonial church in Quito and found a hidden manuscript. The salvaged fragment is a Spanish translation from Quechua of the "Elegy to the Dead of Atahualpa",[112] a poem written by Collahuazo, which describes the sadness and impotence of the Inca people of having lost their king Atahualpa.
+
+Other early Ecuadorian writers include the Jesuits Juan Bautista Aguirre, born in Daule in 1725, and Father Juan de Velasco, born in Riobamba in 1727. De Velasco wrote about the nations and chiefdoms that had existed in the Kingdom of Quito (today Ecuador) before the arrival of the Spanish. His historical accounts are nationalistic, featuring a romantic perspective of precolonial history.
+
+Famous authors from the late colonial and early republic period include Eugenio Espejo, a printer and main author of the first newspaper in Ecuadorian colonial times; Jose Joaquin de Olmedo (born in Guayaquil), famous for his ode to Simón Bolívar titled Victoria de Junin; Juan Montalvo, a prominent essayist and novelist; Juan Leon Mera, famous for his work "Cumanda" or "Tragedy among Savages" and the Ecuadorian National Anthem; Juan A. Martinez with A la Costa';, Dolores Veintimilla;[115] and others.
+
+Contemporary Ecuadorian writers include the novelist Jorge Enrique Adoum; the poet Jorge Carrera Andrade; the essayist Benjamín Carrión; the poets Medardo Angel Silva, Jorge Carrera Andrade, and Luis Alberto Costales; the novelist Enrique Gil Gilbert; the novelist Jorge Icaza (author of the novel Huasipungo, translated to many languages); the short story author Pablo Palacio; and the novelist Alicia Yanez Cossio.
+
+In spite of Ecuador's considerable mystique, it is rarely featured as a setting in contemporary western literature. One exception is "The Ecuadorian Deception," a murder mystery/thriller authored by American Bear Mills. In it, George d'Hout, a website designer from the United States is lured under false pretenses to Guayaquil. A corrupt American archaeologist is behind the plot, believing d'Hout holds the keys to locating a treasure hidden by a buccaneer ancestor. The story is based on a real pirate by the name of George d'Hout who terrorized Guayaquil in the 16th Century.
+
+The best known art styles from Ecuador belonged to the Escuela Quiteña (Quito School), which developed from the 16th to 18th centuries, examples of which are on display in various old churches in Quito. Ecuadorian painters include Eduardo Kingman, Oswaldo Guayasamín, and Camilo Egas from the Indiginist Movement; Manuel Rendon, Jaime Zapata, Enrique Tábara, Aníbal Villacís, Theo Constanté, Luis Molinari, Araceli Gilbert, Judith Gutierrez, Felix Arauz, and Estuardo Maldonado from the Informalist Movement; Teddy Cobeña from expressionism and figurative style[116][117][118] and Luis Burgos Flor with his abstract, futuristic style. The Amerindian people of Tigua, Ecuador, are also world-renowned[citation needed] for their traditional paintings.
+
+Centro de Arte Contemporáneo, Quito
+
+Centro Cultural Metropolitano in the historic center of Quito
+
+The most popular sport in Ecuador, as in most South American countries, is football. Its best known professional teams include; Emelec from Guayaquil also the most popular team in Ecuador, Liga De Quito from Quito; Barcelona S.C. from Guayaquil;  Deportivo Quito, and El Nacional from Quito; Olmedo from Riobamba; and Deportivo Cuenca from Cuenca. Currently the most successful football team in Ecuador is LDU Quito, and it is the only Ecuadorian team that has won the Copa Libertadores, the Copa Sudamericana, and the Recopa Sudamericana; they were also runners-up in the 2008 FIFA Club World Cup. The matches of the Ecuadorian national team are the most-watched sporting events in the country.[citation needed] Ecuador has qualified for the final rounds of the 2002, the 2006, & the 2014 FIFA World Cups. The 2002 FIFA World Cup qualifying campaign was considered a huge success for the country and its inhabitants.[citation needed] The unusually high elevation of the home stadium in Quito often affects the performance of visiting teams. Ecuador finished in 2nd place in the CONMEBOL qualifiers behind Argentina and above the team that would become World Champions, Brazil. In the 2006 FIFA World Cup, Ecuador finished ahead of Poland and Costa Rica finishing second behind Germany in Group A in the 2006 World Cup. They were defeated by England in the second round.
+
+Ecuador has won two medals in the Olympic Games, both gained by 20-km (12 mi) racewalker Jefferson Pérez, who took gold in the 1996 games and silver 12 years later. Pérez also set a world best in the 2003 World Championships of 1:17:21 for the 20-km (12 mi) distance.[119]
+
+In world class professional cycling, Richard Carapaz became the first Ecuadorian to win a Grand Tour. He won the 2019 Giro d'Italia[120]
+
+The current structure of the Ecuadorian public health care system dates back to 1967.[121][122] The Ministry of the Public Health (Ministerio de Salud Pública del Ecuador) is the responsible entity of the regulation and creation of the public health policies and health care plans. The Minister of Public Health is appointed directly by the President of the Republic. The current minister, or Ecuadorian general surgeon, is Margarita Guevara.
+
+The philosophy of the Ministry of Public Health is the social support and service to the most vulnerable population,[123] and its main plan of action lies around communitarian health and preventive medicine.[123]
+
+The public healthcare system allows patients to be treated without an appointment in public general hospitals by general practitioners and specialists in the outpatient clinic (Consulta Externa) at no cost. This is done in the four basic specialties of pediatric, gynecology, clinic medicine, and surgery.[124] There are also public hospitals specialized to treat chronic diseases, target a particular group of the population, or provide better treatment in some medical specialties. Some examples in this group are the Gynecologic Hospitals, or Maternities, Children Hospitals, Geriatric Hospitals, and Oncology Institutes.
+
+Although well-equipped general hospitals are found in the major cities or capitals of provinces, there are basic hospitals in the smaller towns and canton cities for family care consultation and treatments in pediatrics, gynecology, clinical medicine, and surgery.[124]
+
+Community health care centers (Centros de Salud) are found inside metropolitan areas of cities and in rural areas. These are day hospitals that provide treatment to patients whose hospitalization is under 24 hours.[124]
+The doctors assigned to rural communities, where the Amerindian population can be substantial, have small clinics under their responsibility for the treatment of patients in the same fashion as the day hospitals in the major cities. The treatment in this case respects the culture of the community.[124]
+
+The public healthcare system should not be confused with the Ecuadorian Social Security healthcare service, which is dedicated to individuals with formal employment and who are affiliated obligatorily through their employers. Citizens with no formal employment may still contribute to the social security system voluntarily and have access to the medical services rendered by the social security system. The Ecuadorian Institute of Social Security (IESS) has several major hospitals and medical sub-centers under its administration across the nation.[125]
+
+Ecuador currently ranks 20, in most efficient health care countries, compared to 111 back in the year 2000.[126]  Ecuadorians have a life expectancy of 77.1 years.[127] The infant mortality rate is 13 per 1,000 live births,[128] a major improvement from approximately 76 in the early 1980s and 140 in 1950.[129] 23% of children under five are chronically malnourished.[128] Population in some rural areas have no access to potable water, and its supply is provided by mean of water tankers. There are 686 malaria cases per 100,000 people.[130] Basic health care, including doctor's visits, basic surgeries, and basic medications, has been provided free since 2008.[128] However, some public hospitals are in poor condition and often lack necessary supplies to attend the high demand of patients. Private hospitals and clinics are well equipped but still expensive for the majority of the population.
+
+Between 2008 and 2016, new public hospitals have been built, the number of civil servants has increased significantly and salaries have been increased. In 2008, the government introduced universal and compulsory social security coverage. In 2015, corruption remains a problem. Overbilling is recorded in 20% of public establishments and in 80% of private establishments.[131]
+
+The Ecuadorian Constitution requires that all children attend school until they achieve a "basic level of education", which is estimated at nine school years.[133] In 1996, the net primary enrollment rate was 96.9%, and 71.8% of children stayed in school until the fifth grade / age 10.[133] The cost of primary and secondary education is borne by the government, but families often face significant additional expenses such as fees and transportation costs.[133]
+
+Provision of public schools falls far below the levels needed, and class sizes are often very large, and families of limited means often find it necessary to pay for education.[citation needed] In rural areas, only 10% of the children go on to high school.[134] The Ministry of Education states that the mean number of years completed is 6.7.[citation needed]
+
+Ecuador has 61 universities, many of which still confer terminal degrees according to the traditional Spanish education system,[135] honoring a long tradition of having some of the oldest universities in the Americas: University of San Fulgencio, founded in 1586 by the Augustines; San Gregorio Magno University, founded in 1651 by the Jesuits; and University of Santo Tomás of Aquino, founded in 1681 by the Dominican order.
+
+Among the traditional conferred terminal degrees can be noted the doctorate for medicine and law schools or engineering, physics, chemistry, or mathematics for polytechnic or technology institutes. These terminal degrees, as in the case of the PhD in other countries, were the main requirement for an individual to be accepted in academia as a professor or researcher. In the professional realm, a terminal degree granted by an accredited institution automatically provides a professional license to the individual.
+
+However, in 2004, the National Council of Higher Education (CONESUP), started the reorganization of all the degree-granting schemes of the accredited universities in order to pair them with foreign counterparts. The new structure of some careers caused the dropping of subjects, credits, or even the name of the previously conferred diplomas. The terminal degree in law, previously known as JD Juris Doctor (Doctor en Jurisprudencia) was replaced by the one of abogado (attorney) with the exception of the modification of the number of credits to equate it to an undergraduate degree. In the same fashion for medical school, the required time of education was considerably reduced from nine years (the minimum needed to obtain the title of MD in Medicine and Surgery) to almost five, with the provision that the diploma is not terminal anymore, and it is given with the title of médico (medic). Therefore, an MD or PhD in medicine is only to be obtained overseas until the universities adjust themselves to granting schemes and curriculum as in foreign counterparts. Nonetheless, a "médico" can start a career as family practitioner or general medicine physician.
+
+This new reorganization, although very ambitious, lacked the proper path to the homologation of diplomas for highly educated professionals graduated in the country or even for the ones graduated in foreign institutions. One of the points of conflict was the imposition of obtaining foreign degrees to current academicians. As today, a master's degree is a requirement to keep an academic position and at least a foreign PhD to attain or retain the status of rector (president of a university) or décano (dean). For Ecuadorian researchers and many academicians trained in the country, these regulations sounded illogical, disappointing, and unlawful since it appeared a question of a title name conflict rather than specialization or science advancement.
+
+A debate to modify this and other reforms, especially the one which granted control of the Higher Education System by the government, was practically passed with consensus by the multi-partisan National Assembly on August 4, 2010, but vetoed by President Rafael Correa, who wanted to keep the law strictly as it was originally redacted by his political party and SENPLADES (National Secretary of Planning and Development). Due to this change, there are many highly educated professionals and academicians under the old structure but estimated that only 87% of the faculty in public universities have already obtained a master's degree, and fewer than 5% have a PhD (although many of them already have Ecuadorian-granted doctorate degrees).
+
+About 300 institutes of higher education offer two to three years of post-secondary vocational or technical training.
+
+Ecuador is currently placed in 96th position of innovation in technology.[136] The most notable icons in Ecuadorian sciences are the mathematician and cartographer Pedro Vicente Maldonado, born in Riobamba in 1707, and the printer, independence precursor, and medical pioneer Eugenio Espejo, born in 1747 in Quito. Among other notable Ecuadorian scientists and engineers are Lieutenant Jose Rodriguez Labandera,[137] a pioneer who built the first submarine in Latin America in 1837; Reinaldo Espinosa Aguilar (1898–1950), a botanist and biologist of Andean flora; and José Aurelio Dueñas (1880–1961), a chemist and inventor of a method of textile serigraphy.
+
+The major areas of scientific research in Ecuador have been in the medical fields, tropical and infectious diseases treatments, agricultural engineering, pharmaceutical research, and bioengineering. Being a small country and a consumer of foreign technology, Ecuador has favored research supported by entrepreneurship in information technology. The antivirus program Checkprogram, banking protection system MdLock, and Core Banking Software Cobis are products of Ecuadorian development.[138]
+
+The scientific production in hard sciences has been limited due to lack of funding but focused around physics, statistics, and partial differential equations in mathematics.[citation needed] In the case of engineering fields, the majority of scientific production comes from the top three polytechnic institutions: Escuela Superior Politécnica del Litoral - ESPOL, Universidad de Las Fuerzas Armadas - ESPE, and Escuela Politécnica Nacional EPN. The Center for Research and Technology Development in Ecuador is an autonomous center for research and technology development funded by Senecyt.
+
+However, according to Nature, the multidisciplinary scientific journal the top 10 institutions that carry the most outstanding scientific contributions are: Yachay Tech University (Yachay Tech), Escuela Politécnica Nacional (EPN), and Universidad San Francisco de Quito (USFQ).[139]
+
+EPN is known for research and education in the  applied science, astronomy, atmospheric physics, engineering and physical sciences. The Geophysics Institute [140] monitors over the country's volcanoes in the Andes Mountains of Ecuador and in the Galápagos Islands, all of which is part of the Ring of Fire. EPN adopted the polytechnic university model that stresses laboratory instruction in applied science and engineering.
+
+The oldest observatory in South America is the Quito Astronomical Observatory and is located in Quito, Ecuador. The Quito Astronomical Observatory, which gives the global community of a Virtual Telescope System that is connected via the Internet and allows the world to watch by streaming, is managed by EPN.
+
+Contemporary Ecuadorian scientists who have been recognized by international institutions are Eugenia del Pino (born 1945), the first Ecuadorian to be elected to the United States National Academy of Science, and Arturo Villavicencio, who was part of the working group of the IPCC, which shared the 2007 Nobel Peace Prize with Al Gore for their dissemination of the effects of climate change.
+
+Currently, the politics of research and investigation are managed by the National Secretary of Higher Education, Science, and Technology (Senescyt).[141]

en/1796.html.txt

@@ -0,0 +1,83 @@
+
+
+An equinox is commonly regarded as the instant of time when the plane (extended indefinitely in all directions) of Earth's equator passes through the center of the Sun.[3][4] This occurs twice each year, around 20 March and 23 September. In other words, it is the moment at which the center of the visible Sun is directly above the equator.
+
+The word is derived from the Latin aequinoctium, from aequus (equal) and nox (genitive noctis) (night). On the day of an equinox, daytime and nighttime are of approximately equal duration all over the planet. They are not exactly equal, however, due to the angular size of the Sun, atmospheric refraction, and the rapidly changing duration of the length of day that occurs at most latitudes around the equinoxes. Long before conceiving this equality, primitive equatorial cultures noted the day when the Sun rises due east and sets due west, and indeed this happens on the day closest to the astronomically defined event.
+
+In the Northern Hemisphere, the March equinox is called the vernal or spring equinox while the September equinox is called the autumnal or fall equinox. In the Southern Hemisphere, the reverse is true. The dates slightly vary due to leap years and other factors.[5]
+
+Since the Moon (and to a lesser extent the planets) causes Earth's orbit to slightly vary from a perfect ellipse, the equinox is officially defined by the Sun's more regular ecliptic longitude rather than by its declination. The instants of the equinoxes are currently defined to be when the apparent geocentric longitude of the Sun is 0° and 180°.[6]
+
+Systematically observing the sunrise, people discovered that it occurs between two extreme locations at the horizon and eventually noted the midpoint between the two. Later it was realized that this happens on a day when the durations of the day and the night are practically equal and the word "equinox" comes from Latin aequus, meaning "equal",  and nox, meaning "night".
+
+In the northern hemisphere, the vernal equinox (March) conventionally marks the beginning of spring in most cultures and is considered the start of the New Year in the Assyrian calendar, Hindu, and the Persian or Iranian calendars,[a] while the autumnal equinox (September) marks the beginning of autumn.[7]
+
+Illumination of Earth by the Sun at the equinox
+
+The relation between the Earth, Sun, and stars at the March equinox. From Earth's perspective, the Sun appears to move along the ecliptic (red), which is tilted compared to the celestial equator (white).
+
+Diagram of the Earth's seasons as seen from the north. Far right: December solstice.
+
+Diagram of the Earth's seasons as seen from the south. Far left: June solstice.
+
+The equinoxes are the only times when the solar terminator (the "edge" between night and day) is perpendicular to the equator. As a result, the northern and southern hemispheres are equally illuminated.
+
+For the same reason, this is also the time when the Sun rises for an observer at one of Earth's rotational poles and sets at the other; for a brief period, both North and South Poles are in daylight.[b]
+
+In other words, the equinoxes are the only times when the subsolar point is on the equator, meaning that the Sun is exactly overhead at a point on the equatorial line. The subsolar point crosses the equator moving northward at the March equinox and southward at the September equinox.
+
+When Julius Caesar established the Julian calendar in 45 BC, he set 25 March as the date of the spring equinox; this was already the starting day of the year in the Persian and Indian calendars. Because the Julian year is longer than the tropical year by about 11.3 minutes on average (or 1 day in 128 years), the calendar "drifted" with respect to the two equinoxes – so that in 300 AD the spring equinox occurred on about 21 March, and by 1500 AD it had drifted backwards to 11 March.[citation needed]
+
+This drift induced Pope Gregory XIII to establish the modern Gregorian calendar. The Pope wanted to continue to conform with the edicts of the Council of Nicaea in 325 AD concerning the date of Easter, which means he wanted to move the vernal equinox to the date on which it fell at that time (21 March is the day allocated to it in the Easter table of the Julian calendar), and to maintain it at around that date in the future, which he achieved by reducing the number of leap years from 100 to 97 every 400 years. However, there remained a small residual variation in the date and time of the vernal equinox of about ±27 hours from its mean position, virtually all because the distribution of 24 hour centurial leap-days causes large jumps (see Gregorian calendar leap solstice). This in turn raised the possibility that it could fall on 22 March, and thus Easter Day might theoretically commence before the equinox. The consulting astronomers chose the appropriate number of days to omit so that the equinox would swing from 19 to 21 March but never fall on 22 March (within Europe).
+
+The dates of the equinoxes change progressively during the leap-year cycle, because the Gregorian calendar year is not commensurate with the period of the Earth's revolution about the Sun. It is only after a complete Gregorian leap-year cycle of 400 years that the seasons commence at approximately the same time. In the 21st century the earliest March equinox will be 19 March 2096, while the latest was 21 March 2003. The earliest September equinox will be 21 September 2096 while the latest was 23 September 2003 (Universal Time).[5]
+
+Day is usually defined as the period when sunlight reaches the ground in the absence of local obstacles.[citation needed] On the date of the equinox, the center of the Sun spends a roughly equal amount of time above and below the horizon at every location on the Earth, so night and day are about the same length.  Sunrise and sunset can be defined in several ways, but a widespread definition is the time that the top limb of the Sun is level with the horizon.[17] With this definition, the day is longer than the night at the equinoxes:[3]
+
+In sunrise/sunset tables, the atmospheric refraction is assumed to be 34 arcminutes, and the assumed semidiameter (apparent radius) of the Sun is 16 arcminutes. (The apparent radius varies slightly depending on time of year, slightly larger at perihelion in January than aphelion in July, but the difference is comparatively small.) Their combination means that when the upper limb of the Sun is on the visible horizon, its centre is 50 arcminutes below the geometric horizon, which is the intersection with the celestial sphere of a horizontal plane through the eye of the observer.[18]
+
+These effects make the day about 14 minutes longer than the night at the equator and longer still towards the poles. The real equality of day and night only happens in places far enough from the equator to have a seasonal difference in day length of at least 7 minutes,[19] actually occurring a few days towards the winter side of each equinox.
+
+The times of sunset and sunrise vary with the observer's location (longitude and latitude), so the dates when day and night are equal also depend upon the observer's location.
+
+A third correction for the visual observation of a sunrise (or sunset) is the angle between the apparent horizon as seen by an observer and the geometric (or sensible) horizon. This is known as the dip of the horizon and varies from 3 arcminutes for a viewer standing on the sea shore to 160 arcminutes for a mountaineer on Everest.[20] The effect of a larger dip on taller objects (reaching over 2½° of arc on Everest) accounts for the phenomenon of snow on a mountain peak turning gold in the sunlight long before the lower slopes are illuminated.
+
+The date on which the day and night are exactly the same is known as an equilux; the neologism, believed to have been coined in the 1980s, achieved more widespread recognition in the 21st century.[c] At the most precise measurements, a true equilux is rare, because the lengths of day and night change more rapidly than any other time of the year around the equinoxes. In the mid-latitudes, daylight increases or decreases by about three minutes per day at the equinoxes, and thus adjacent days and nights only reach within one minute of each other. The date of the closest approximation of the equilux varies slightly by latitude; in the mid-latitudes, it occurs a few days before the spring equinox and after the fall equinox in each respective hemisphere.
+
+In the half-year centered on the June solstice, the Sun rises north of east and sets north of west, which means longer days with shorter nights for the northern hemisphere and shorter days with longer nights for the southern hemisphere. In the half-year centered on the December solstice, the Sun rises south of east and sets south of west and the durations of day and night are reversed.
+
+Also on the day of an equinox, the Sun rises everywhere on Earth (except at the poles) at about 06:00 and sets at about 18:00 (local solar time). These times are not exact for several reasons:
+
+Some of the statements above can be made clearer by picturing the day arc (i.e., the path along which the Sun appears to move across the sky). The pictures show this for every hour on equinox day. In addition, some 'ghost' suns are also indicated below the horizon, up to 18° below it; the Sun in such areas still causes twilight. The depictions presented below can be used for both the northern and the southern hemispheres. The observer is understood to be sitting near the tree on the island depicted in the middle of the ocean; the green arrows give cardinal directions.
+
+The following special cases are depicted:
+
+Day arc at 0° latitude (equator)The arc passes through the zenith, resulting in any purely vertical object (such as an obelisk or pillar) having no shadow at high noon.
+
+Day arc at 20° latitudeThe Sun culminates at 70° altitude and its path at sunrise and sunset occurs at a steep 70° angle to the horizon. Twilight still lasts about one hour.
+
+Day arc at 50° latitudeTwilight lasts almost two hours.
+
+Day arc at 70° latitudeThe Sun culminates at no more than 20° altitude and its daily path at sunrise and sunset is at a shallow 20° angle to the horizon. Twilight lasts for more than four hours.
+
+Day arc at 90° latitude (pole)If it were not for atmospheric refraction, the Sun would be on the horizon all the time.
+
+The March equinox occurs about when the Sun appears to cross the celestial equator northward. In the Northern Hemisphere, the term vernal point is used for the time of this occurrence and for the precise direction in space where the Sun exists at that time. This point is the origin of some celestial coordinate systems, which are usually rooted to an astronomical epoch since it gradually varies (precesses) over time:
+
+Strictly speaking, at the equinox, the Sun's ecliptic longitude is zero. Its latitude will not be exactly zero, since Earth is not exactly in the plane of the ecliptic. Its declination will not be exactly zero either. The mean ecliptic is defined by the barycenter of Earth and the Moon combined, so the Earth wanders slightly above and below the ecliptic due to the orbital tilt of the Moon.[26] The modern definition of equinox is the instants when the Sun's apparent geocentric longitude is 0° (northward equinox) or 180° (southward equinox).[27][28][29] See the adjacent diagram.
+
+Because of the precession of the Earth's axis, the position of the vernal point on the celestial sphere changes over time, and the equatorial and the ecliptic coordinate systems change accordingly. Thus when specifying celestial coordinates for an object, one has to specify at what time the vernal point and the celestial equator are taken. That reference time is called the equinox of date.[30]
+
+The upper culmination of the vernal point is considered the start of the sidereal day for the observer. The hour angle of the vernal point is, by definition, the observer's sidereal time.
+
+Using the current official IAU constellation boundaries – and taking into account the variable precession speed and the rotation of the celestial equator – the equinoxes shift through the constellations as follows[31] (expressed in astronomical year numbering when the year 0 = 1 BC, −1 = 2 BC, etc.):
+
+The equinoxes are sometimes regarded as the start of spring and autumn. A number of traditional harvest festivals  are celebrated on the date of the equinoxes.
+
+One effect of equinoctial periods is the temporary disruption of communications satellites. For all geostationary satellites, there are a few days around the equinox when the Sun goes directly behind the satellite relative to Earth (i.e. within the beam-width of the ground-station antenna) for a short period each day. The Sun's immense power and broad radiation spectrum overload the Earth station's reception circuits with noise and, depending on antenna size and other factors, temporarily disrupt or degrade the circuit. The duration of those effects varies but can range from a few minutes to an hour. (For a given frequency band, a larger antenna has a narrower beam-width and hence experiences shorter duration "Sun outage" windows.)[32]
+
+Satellites in geostationary orbit also experience difficulties maintaining power during the equinox, due to the fact that they now have to travel through Earth's shadow and rely only on battery power. Usually, a satellite will travel either north or south of the Earth's shadow due to its shifted axis throughout the year. During the equinox, since geostationary satellites are situated above the Equator, they will be put into Earth's shadow for the longest duration all year.[33]
+
+Equinoxes occur on any planet with a tilted rotational axis. A dramatic example is Saturn, where the equinox places its ring system edge-on facing the Sun. As a result, they are visible only as a thin line when seen from Earth. When seen from above – a view seen during an equinox for the first time from the Cassini space probe in 2009 – they receive very little sunshine, indeed more planetshine than light from the Sun.[34] This phenomenon occurs once every 14.7 years on average, and can last a few weeks before and after the exact equinox. Saturn's most recent equinox was on 11 August 2009, and its next will take place on 6 May 2025.[35]
+
+Mars's most recent equinox was on 8 April 2020 (northern autumn), and the next will be on 7 February 2021 (northern spring).[36]

en/1797.html.txt

@@ -0,0 +1,83 @@
+
+
+An equinox is commonly regarded as the instant of time when the plane (extended indefinitely in all directions) of Earth's equator passes through the center of the Sun.[3][4] This occurs twice each year, around 20 March and 23 September. In other words, it is the moment at which the center of the visible Sun is directly above the equator.
+
+The word is derived from the Latin aequinoctium, from aequus (equal) and nox (genitive noctis) (night). On the day of an equinox, daytime and nighttime are of approximately equal duration all over the planet. They are not exactly equal, however, due to the angular size of the Sun, atmospheric refraction, and the rapidly changing duration of the length of day that occurs at most latitudes around the equinoxes. Long before conceiving this equality, primitive equatorial cultures noted the day when the Sun rises due east and sets due west, and indeed this happens on the day closest to the astronomically defined event.
+
+In the Northern Hemisphere, the March equinox is called the vernal or spring equinox while the September equinox is called the autumnal or fall equinox. In the Southern Hemisphere, the reverse is true. The dates slightly vary due to leap years and other factors.[5]
+
+Since the Moon (and to a lesser extent the planets) causes Earth's orbit to slightly vary from a perfect ellipse, the equinox is officially defined by the Sun's more regular ecliptic longitude rather than by its declination. The instants of the equinoxes are currently defined to be when the apparent geocentric longitude of the Sun is 0° and 180°.[6]
+
+Systematically observing the sunrise, people discovered that it occurs between two extreme locations at the horizon and eventually noted the midpoint between the two. Later it was realized that this happens on a day when the durations of the day and the night are practically equal and the word "equinox" comes from Latin aequus, meaning "equal",  and nox, meaning "night".
+
+In the northern hemisphere, the vernal equinox (March) conventionally marks the beginning of spring in most cultures and is considered the start of the New Year in the Assyrian calendar, Hindu, and the Persian or Iranian calendars,[a] while the autumnal equinox (September) marks the beginning of autumn.[7]
+
+Illumination of Earth by the Sun at the equinox
+
+The relation between the Earth, Sun, and stars at the March equinox. From Earth's perspective, the Sun appears to move along the ecliptic (red), which is tilted compared to the celestial equator (white).
+
+Diagram of the Earth's seasons as seen from the north. Far right: December solstice.
+
+Diagram of the Earth's seasons as seen from the south. Far left: June solstice.
+
+The equinoxes are the only times when the solar terminator (the "edge" between night and day) is perpendicular to the equator. As a result, the northern and southern hemispheres are equally illuminated.
+
+For the same reason, this is also the time when the Sun rises for an observer at one of Earth's rotational poles and sets at the other; for a brief period, both North and South Poles are in daylight.[b]
+
+In other words, the equinoxes are the only times when the subsolar point is on the equator, meaning that the Sun is exactly overhead at a point on the equatorial line. The subsolar point crosses the equator moving northward at the March equinox and southward at the September equinox.
+
+When Julius Caesar established the Julian calendar in 45 BC, he set 25 March as the date of the spring equinox; this was already the starting day of the year in the Persian and Indian calendars. Because the Julian year is longer than the tropical year by about 11.3 minutes on average (or 1 day in 128 years), the calendar "drifted" with respect to the two equinoxes – so that in 300 AD the spring equinox occurred on about 21 March, and by 1500 AD it had drifted backwards to 11 March.[citation needed]
+
+This drift induced Pope Gregory XIII to establish the modern Gregorian calendar. The Pope wanted to continue to conform with the edicts of the Council of Nicaea in 325 AD concerning the date of Easter, which means he wanted to move the vernal equinox to the date on which it fell at that time (21 March is the day allocated to it in the Easter table of the Julian calendar), and to maintain it at around that date in the future, which he achieved by reducing the number of leap years from 100 to 97 every 400 years. However, there remained a small residual variation in the date and time of the vernal equinox of about ±27 hours from its mean position, virtually all because the distribution of 24 hour centurial leap-days causes large jumps (see Gregorian calendar leap solstice). This in turn raised the possibility that it could fall on 22 March, and thus Easter Day might theoretically commence before the equinox. The consulting astronomers chose the appropriate number of days to omit so that the equinox would swing from 19 to 21 March but never fall on 22 March (within Europe).
+
+The dates of the equinoxes change progressively during the leap-year cycle, because the Gregorian calendar year is not commensurate with the period of the Earth's revolution about the Sun. It is only after a complete Gregorian leap-year cycle of 400 years that the seasons commence at approximately the same time. In the 21st century the earliest March equinox will be 19 March 2096, while the latest was 21 March 2003. The earliest September equinox will be 21 September 2096 while the latest was 23 September 2003 (Universal Time).[5]
+
+Day is usually defined as the period when sunlight reaches the ground in the absence of local obstacles.[citation needed] On the date of the equinox, the center of the Sun spends a roughly equal amount of time above and below the horizon at every location on the Earth, so night and day are about the same length.  Sunrise and sunset can be defined in several ways, but a widespread definition is the time that the top limb of the Sun is level with the horizon.[17] With this definition, the day is longer than the night at the equinoxes:[3]
+
+In sunrise/sunset tables, the atmospheric refraction is assumed to be 34 arcminutes, and the assumed semidiameter (apparent radius) of the Sun is 16 arcminutes. (The apparent radius varies slightly depending on time of year, slightly larger at perihelion in January than aphelion in July, but the difference is comparatively small.) Their combination means that when the upper limb of the Sun is on the visible horizon, its centre is 50 arcminutes below the geometric horizon, which is the intersection with the celestial sphere of a horizontal plane through the eye of the observer.[18]
+
+These effects make the day about 14 minutes longer than the night at the equator and longer still towards the poles. The real equality of day and night only happens in places far enough from the equator to have a seasonal difference in day length of at least 7 minutes,[19] actually occurring a few days towards the winter side of each equinox.
+
+The times of sunset and sunrise vary with the observer's location (longitude and latitude), so the dates when day and night are equal also depend upon the observer's location.
+
+A third correction for the visual observation of a sunrise (or sunset) is the angle between the apparent horizon as seen by an observer and the geometric (or sensible) horizon. This is known as the dip of the horizon and varies from 3 arcminutes for a viewer standing on the sea shore to 160 arcminutes for a mountaineer on Everest.[20] The effect of a larger dip on taller objects (reaching over 2½° of arc on Everest) accounts for the phenomenon of snow on a mountain peak turning gold in the sunlight long before the lower slopes are illuminated.
+
+The date on which the day and night are exactly the same is known as an equilux; the neologism, believed to have been coined in the 1980s, achieved more widespread recognition in the 21st century.[c] At the most precise measurements, a true equilux is rare, because the lengths of day and night change more rapidly than any other time of the year around the equinoxes. In the mid-latitudes, daylight increases or decreases by about three minutes per day at the equinoxes, and thus adjacent days and nights only reach within one minute of each other. The date of the closest approximation of the equilux varies slightly by latitude; in the mid-latitudes, it occurs a few days before the spring equinox and after the fall equinox in each respective hemisphere.
+
+In the half-year centered on the June solstice, the Sun rises north of east and sets north of west, which means longer days with shorter nights for the northern hemisphere and shorter days with longer nights for the southern hemisphere. In the half-year centered on the December solstice, the Sun rises south of east and sets south of west and the durations of day and night are reversed.
+
+Also on the day of an equinox, the Sun rises everywhere on Earth (except at the poles) at about 06:00 and sets at about 18:00 (local solar time). These times are not exact for several reasons:
+
+Some of the statements above can be made clearer by picturing the day arc (i.e., the path along which the Sun appears to move across the sky). The pictures show this for every hour on equinox day. In addition, some 'ghost' suns are also indicated below the horizon, up to 18° below it; the Sun in such areas still causes twilight. The depictions presented below can be used for both the northern and the southern hemispheres. The observer is understood to be sitting near the tree on the island depicted in the middle of the ocean; the green arrows give cardinal directions.
+
+The following special cases are depicted:
+
+Day arc at 0° latitude (equator)The arc passes through the zenith, resulting in any purely vertical object (such as an obelisk or pillar) having no shadow at high noon.
+
+Day arc at 20° latitudeThe Sun culminates at 70° altitude and its path at sunrise and sunset occurs at a steep 70° angle to the horizon. Twilight still lasts about one hour.
+
+Day arc at 50° latitudeTwilight lasts almost two hours.
+
+Day arc at 70° latitudeThe Sun culminates at no more than 20° altitude and its daily path at sunrise and sunset is at a shallow 20° angle to the horizon. Twilight lasts for more than four hours.
+
+Day arc at 90° latitude (pole)If it were not for atmospheric refraction, the Sun would be on the horizon all the time.
+
+The March equinox occurs about when the Sun appears to cross the celestial equator northward. In the Northern Hemisphere, the term vernal point is used for the time of this occurrence and for the precise direction in space where the Sun exists at that time. This point is the origin of some celestial coordinate systems, which are usually rooted to an astronomical epoch since it gradually varies (precesses) over time:
+
+Strictly speaking, at the equinox, the Sun's ecliptic longitude is zero. Its latitude will not be exactly zero, since Earth is not exactly in the plane of the ecliptic. Its declination will not be exactly zero either. The mean ecliptic is defined by the barycenter of Earth and the Moon combined, so the Earth wanders slightly above and below the ecliptic due to the orbital tilt of the Moon.[26] The modern definition of equinox is the instants when the Sun's apparent geocentric longitude is 0° (northward equinox) or 180° (southward equinox).[27][28][29] See the adjacent diagram.
+
+Because of the precession of the Earth's axis, the position of the vernal point on the celestial sphere changes over time, and the equatorial and the ecliptic coordinate systems change accordingly. Thus when specifying celestial coordinates for an object, one has to specify at what time the vernal point and the celestial equator are taken. That reference time is called the equinox of date.[30]
+
+The upper culmination of the vernal point is considered the start of the sidereal day for the observer. The hour angle of the vernal point is, by definition, the observer's sidereal time.
+
+Using the current official IAU constellation boundaries – and taking into account the variable precession speed and the rotation of the celestial equator – the equinoxes shift through the constellations as follows[31] (expressed in astronomical year numbering when the year 0 = 1 BC, −1 = 2 BC, etc.):
+
+The equinoxes are sometimes regarded as the start of spring and autumn. A number of traditional harvest festivals  are celebrated on the date of the equinoxes.
+
+One effect of equinoctial periods is the temporary disruption of communications satellites. For all geostationary satellites, there are a few days around the equinox when the Sun goes directly behind the satellite relative to Earth (i.e. within the beam-width of the ground-station antenna) for a short period each day. The Sun's immense power and broad radiation spectrum overload the Earth station's reception circuits with noise and, depending on antenna size and other factors, temporarily disrupt or degrade the circuit. The duration of those effects varies but can range from a few minutes to an hour. (For a given frequency band, a larger antenna has a narrower beam-width and hence experiences shorter duration "Sun outage" windows.)[32]
+
+Satellites in geostationary orbit also experience difficulties maintaining power during the equinox, due to the fact that they now have to travel through Earth's shadow and rely only on battery power. Usually, a satellite will travel either north or south of the Earth's shadow due to its shifted axis throughout the year. During the equinox, since geostationary satellites are situated above the Equator, they will be put into Earth's shadow for the longest duration all year.[33]
+
+Equinoxes occur on any planet with a tilted rotational axis. A dramatic example is Saturn, where the equinox places its ring system edge-on facing the Sun. As a result, they are visible only as a thin line when seen from Earth. When seen from above – a view seen during an equinox for the first time from the Cassini space probe in 2009 – they receive very little sunshine, indeed more planetshine than light from the Sun.[34] This phenomenon occurs once every 14.7 years on average, and can last a few weeks before and after the exact equinox. Saturn's most recent equinox was on 11 August 2009, and its next will take place on 6 May 2025.[35]
+
+Mars's most recent equinox was on 8 April 2020 (northern autumn), and the next will be on 7 February 2021 (northern spring).[36]

en/1798.html.txt

@@ -0,0 +1,179 @@
+
+
+
+
+The Germany national football team (German: Deutsche Fußballnationalmannschaft or Die Mannschaft) represents Germany in men's international football and played its first match in 1908.[5] The team is governed by the German Football Association (Deutscher Fußball-Bund), founded in 1900.[9][10] Between 1949 and 1990, separate German national teams were recognised by FIFA due to Allied occupation and division: the DFB's team representing the Federal Republic of Germany (named West Germany from 1949–1990), the Saarland team representing the Saar Protectorate (1950–1956) and the East German team representing the German Democratic Republic (1952–1990). The latter two were absorbed along with their records;[11][12] the present team represents the reunified Federal Republic. The official name and code "Germany FR (FRG)" was shortened to "Germany (GER)" following reunification in 1990.
+
+Germany is one of the most successful national teams in international competitions, having won four World Cups (1954, 1974, 1990, 2014), three European Championships (1972, 1980, 1996), and one Confederations Cup (2017).[9] They have also been runners-up three times in the European Championships, four times in the World Cup, and a further four third-place finishes at World Cups.[9] East Germany won Olympic Gold in 1976.[13]
+
+Germany is the only nation to have won both the FIFA World Cup and the FIFA Women's World Cup.[14][15]
+
+At the end of the 2014 World Cup, Germany earned the highest Elo rating of any national football team in history, with a record 2,205 points.[16] Germany is also the only European nation that has won a FIFA World Cup in the Americas. The manager of the national team is Joachim Löw.
+
+Between 1899 and 1901, prior to the formation of a national team, there were five unofficial international matches between German and English selection teams, which all ended as large defeats for the German teams. Eight years after the establishment of the German Football Association (DFB), the first official match of the Germany national football team[17] was played on 5 April 1908, against Switzerland in Basel, with the Swiss winning 5–3.[5]
+
+Gottfried Fuchs scored a world record 10 goals for Germany in a 16–0 win against Russia at the 1912 Olympics in Stockholm on 1 July, becoming the top scorer of the tournament; his international record was not surpassed until 2001 when Australia's Archie Thompson scored 13 goals in a 31–0 defeat of American Samoa.[18] He was Jewish, and the German Football Association erased all references to him from their records between 1933 and 1945.[19][20] As of 2016, he was still the top German scorer for one match.[21]
+
+The first match after World War I in 1920, the first match after World War II in 1950 when Germany was still banned from most international competitions, and the first match in 1990 with former East German players were all against Switzerland as well. Germany's first championship title was even won in Switzerland in 1954.
+
+At that time the players were selected by the DFB, as there was no dedicated coach. The first manager of the Germany national team was Otto Nerz, a school teacher from Mannheim, who served in the role from 1926 to 1936.[22] The German FA could not afford travel to Uruguay for the first World Cup staged in 1930 during the Great Depression, but finished third in the 1934 World Cup in their first appearance in the competition. After a poor showing at the 1936 Olympic Games in Berlin, Sepp Herberger became coach. In 1937 he put together a squad which was soon nicknamed the Breslau Elf (the Breslau Eleven) in recognition of their 8–0 win over Denmark in the then German city of Breslau, Lower Silesia (now Wrocław, Poland).[23][24]
+
+After Austria became part of Germany in the Anschluss of March 1938, that country's national team – one of Europe's best sides at the time due to professionalism – was disbanded despite having already qualified for the 1938 World Cup. Nazi politicians ordered five or six ex-Austrian players, from the clubs Rapid Vienna, Austria Vienna, and First Vienna FC, to join the all-German team on short notice in a staged show of unity for political reasons. In the 1938 World Cup that began on 4 June, this "united" German team managed only a 1–1 draw against Switzerland and then lost the replay 2–4 in front of a hostile crowd in Paris, France. That early exit stands as Germany's worst World Cup result, and one of just two occasions the team failed to progress the group stage (the next would not occur until the 2018 tournament).
+
+During World War II, the team played over 30 international games between September 1939 and November 1942. National team games were then suspended, as most players had to join the armed forces. Many of the national team players were gathered together under coach Herberger as Rote Jäger through the efforts of a sympathetic air force officer trying to protect the footballers from the most dangerous wartime service.
+
+After World War II, Germany was banned from competition in most sports until 1950. The DFB was not a full member of FIFA, and none of the three new German states – West Germany, East Germany, and Saarland – entered the 1950 World Cup qualifiers.
+
+The Federal Republic of Germany, which was referred to as West Germany, continued the DFB. With recognition by FIFA and UEFA, the DFB maintained and continued the record of the pre-war team. Switzerland was once again the first team that played West Germany in 1950.[25]
+West Germany qualified for the 1954 World Cup.
+
+The Saarland, under French control between 1947 and 1956, did not join French organisations, and was barred from participating in pan-German ones. It sent their own team to the 1952 Summer Olympics and to the 1954 World Cup qualifiers. In 1957, Saarland acceded to the Federal Republic of Germany.
+
+In 1949, the communist German Democratic Republic (East Germany) was founded. In 1952 the Deutscher Fußball-Verband der DDR (DFV) was established and the East Germany national football team took to the field. They were the only team to beat the 1974 FIFA World Cup winning West Germans in the only meeting of the two sides of the divided nation. East Germany won the gold medal at the 1976 Olympics. After German reunification in 1990, the eastern football competition was reintegrated into the DFB.
+
+West Germany, captained by Fritz Walter, met in the 1954 World Cup against Turkey, Yugoslavia and Austria. When playing favourites Hungary in the group stage, Germany lost 3–8. West Germany met the Hungarian "Mighty Magyars" again in the final. Hungary had gone unbeaten for 32 consecutive matches. In an upset, West Germany won 3–2, with Helmut Rahn scoring the winning goal.[26] The success is called "The Miracle of Bern" (Das Wunder von Bern).[27]
+
+After finishing fourth in the 1958 World Cup and reaching only the quarter-finals in the 1962 World Cup, the DFB made changes. Professionalism was introduced, and the best clubs from the various Regionalligas were assembled into the new Bundesliga. In 1964, Helmut Schön took over as coach, replacing Herberger who had been in office for 28 years.
+
+In the 1966 World Cup, West Germany reached the final after beating the USSR in the semi-final, facing hosts England. In extra time, the first goal by Geoff Hurst was one of the most contentious goals in the history of the World Cup: the linesman signalled the ball had crossed the line for a goal, after bouncing down from the crossbar, when replays showed it did not appear to have fully crossed the line. Hurst then scored another goal giving England a 4–2 win.[28][29]
+
+West Germany in the 1970 World Cup knocked England out in the quarter-finals 3–2, before they suffered a 4–3 extra-time loss in the semi-final against Italy. This match with five goals in extra time is one of the most dramatic in World Cup history, and is called the "Game of the Century" in both Italy and Germany.[30][31] West Germany claimed third by beating Uruguay 1–0. Gerd Müller finished as the tournament's top scorer with 10 goals.
+
+In 1971, Franz Beckenbauer became captain of the national team, and he led West Germany to victory at the European Championship at Euro 1972, defeating the Soviet Union 3–0 in the final.[32][33]
+
+As hosts of the 1974 World Cup, they won their second World Cup, defeating the Netherlands 2–1 in the final in Munich.[34]
+Two matches in the 1974 World Cup stood out for West Germany. The first group stage saw a politically charged match as West Germany played a game against East Germany. The East Germans won 1–0.[35] The West Germans advanced to the final against the Johan Cruijff-led Dutch team and their brand of "Total Football". The Dutch took the lead from a penalty. However, West Germany tied the match on a penalty by Paul Breitner, and won it with Gerd Müller's fine finish soon after.[36][37]
+
+West Germany failed to defend their titles in the next two major international tournaments. They lost to Czechoslovakia in the final of Euro 1976 in a penalty shootout 5–3.[38] Since that loss, Germany has not lost a penalty shootout in major international tournaments.[39]
+
+In the 1978 World Cup, Germany was eliminated in the second group stage after losing 3–2 to Austria. Schön retired as coach afterward, and the post was taken over by his assistant, Jupp Derwall.
+
+West Germany's first tournament under Derwall was successful, as they earned their second European title at Euro 1980 after defeating Belgium 2–1 in the final.[40] West Germany reached the final of the 1982 World Cup, though not without difficulties. They were upset 1–2 by Algeria in their first match,[41] but advanced to the second round with a controversial 1–0 win over Austria. In the semi-final against France, they tied the match 3–3 and won the penalty shootout 5–4.[42][43] In the final, they were defeated by Italy 1–3.[44]
+
+During this period, West Germany's Gerd Müller racked up fourteen goals in two World Cups (1970 and 1974). His ten goals in 1970 are the third-most ever in a tournament. (Müller's all-time World Cup record of 14 goals was broken by Ronaldo in 2006; this was then further broken by Miroslav Klose in 2014 with 16 goals).[45]
+
+After West Germany were eliminated in the first round of Euro 1984, Franz Beckenbauer returned to the national team to replace Derwall as coach.[46] At the 1986 World Cup in Mexico, West Germany finished as runners-up for the second consecutive tournament after beating France 2–0 in the semi-finals, but losing to the Diego Maradona-led Argentina in the final, 2–3.[47][48] In Euro 1988, West Germany's hopes of winning the tournament on home soil were spoiled by the Netherlands, as the Dutch beat them 2–1 in the semi-finals.[49]
+
+At the 1990 World Cup in Italy, West Germany won their third World Cup title, in its unprecedented third consecutive final appearance.[50] Captained by Lothar Matthäus, they defeated Yugoslavia (4–1), UAE (5–1), the Netherlands (2–1), Czechoslovakia (1–0), and England (1–1, 4–3 on penalty kicks) on the way to a final rematch against Argentina, played in the Italian capital of Rome.[51][52] West Germany won 1–0, with the only goal being a penalty scored in the 85th minute by Andreas Brehme.[50] Beckenbauer, who won the World Cup as the national team's captain in 1974, thus became the first person to win the World Cup as both captain and coach.[46]
+
+Prior to 1984, Olympic football was an amateur event, meaning that only non-professional players could participate. Due to this, West Germany was never able to achieve the same degree of success at the Olympics as at the World Cup, with the first medal coming in the 1988 Olympics, when they won the bronze medal. It took Germany 28 years to participate at the Olympics again in 2016, this time reaching the final and winning a silver medal. West Germany also reached the second round in both 1972 and 1984. On the other hand, East Germany did far better, winning a gold, a silver and two bronze medals (one representing the United Team of Germany).
+
+In February 1990, months after the fall of the Berlin Wall, East Germany and West Germany were drawn together in UEFA Euro 1992 qualifying Group 5. In November 1990, the East German association Deutscher Fußball-Verband integrated into the DFB, by which time the East German team had ceased operations, playing its last match on 12 September 1990. The unified German national team completed the European Championship qualifying group. The East German 1990–91 league continued, with a restructuring of German leagues in 1991–92. The first game with a unified German team was against Sweden on 10 October.[53]
+
+After the 1990 World Cup, assistant Berti Vogts took over as the national team coach from the retiring Beckenbauer. In Euro 1992, Germany reached the final, but lost 0–2 to underdogs Denmark.[54]
+In the 1994 World Cup, they were upset 1–2 in the quarterfinals by Bulgaria.[55][56]
+
+Reunified Germany won its first major international title at Euro 1996, becoming European champions for the third time.[57] They defeated hosts England in the semi-finals,[58] and the Czech Republic 2–1 in the final on a golden goal in extra time.[59]
+
+However, in the 1998 World Cup, Germany were eliminated in the quarterfinals in a 0–3 defeat to Croatia, all goals being scored after defender Christian Wörns received a straight red card.[60] Vogts stepped down afterwards and was replaced by Erich Ribbeck.[61]
+
+In Euro 2000, the team went out in the first round, drawing with Romania, then suffering a 1–0 defeat to England and were routed 3–0 by Portugal (which fielded their backup players, having already advanced).[62] Ribbeck resigned, and was replaced by Rudi Völler.[63]
+
+Coming into the 2002 World Cup, expectations of the German team were low due to poor results in the qualifiers and not directly qualifying for the finals for the first time. The team advanced through group play, and in the knockout stages they produced three consecutive 1–0 wins against Paraguay,[64] the United States,[65] and co-hosts South Korea. Oliver Neuville scored two minutes from time against Paraguay and Michael Ballack scored both goals in the US and South Korea games, although he picked up a second yellow card against South Korea for a tactical foul and was suspended for the subsequent match.[66] This set up a final against Brazil, the first World Cup meeting between the two. Germany lost 0–2 thanks to two Ronaldo goals.[67] Nevertheless, German captain and goalkeeper Oliver Kahn won the Golden Ball,[68] the first time in the World Cup that a goalkeeper was named the best player of the tournament.[69]
+
+Germany again exited in the first round of Euro 2004, drawing their first two matches and losing the third to the Czech Republic (who had fielded a second-string team).[70] Völler resigned afterwards, and Jürgen Klinsmann was appointed head coach.[71][72]
+
+Klinsmann's main task was to lead the national team to a good showing at the 2006 World Cup in Germany. Klinsmann relieved goalkeeper Kahn of the captaincy and announced that Kahn and longtime backup Jens Lehmann would be competing for the position of starting goaltender, a decision that angered Kahn and Lehmann eventually won that contest.[73] Expectations for the team were low, which was not helped by veteran defender Christian Wörns being dropped (after Wörns criticised Klinsmann for designating him only as a backup player on the squad), a choice roundly panned in Germany. Italy routed Germany 4–1 in a March exhibition game, and Klinsmann bore the brunt of the criticism as the team was ranked only 22nd in the world entering the 2006 FIFA World Cup.[74]
+
+As World Cup hosts, Germany won all three group-stage matches to finish top of their group. The team defeated Sweden 2–0 in the round of 16.[75]
+Germany faced Argentina in the quarter-finals. The match ended 1–1, and Germany won the penalty shootout 4–2.[76]
+In the semi-final against Italy, the match was scoreless until near the end of extra time when Germany conceded two goals.[77]
+In the third place match, Germany defeated Portugal 3–1.[78]
+Miroslav Klose was awarded the Golden Boot for his tournament-leading five goals.[79]
+
+Germany's entry into the Euro 2008 qualifying round was marked by the promotion of Joachim Löw to head coach, since Klinsmann resigned.[80]
+At UEFA Euro 2008, Germany won two out of three matches in group play to advance to the knockout round.[81]
+They defeated Portugal 3–2 in the quarterfinal,[82]
+and won their semi-final against Turkey.[83]
+Germany lost the final against Spain 0–1, finishing as the runners-up.[84]
+
+In the 2010 World Cup, Germany won the group and advanced to the knockout stage. In the round of 16, Germany defeated England 4–1.[85] The game controversially had a valid goal by Frank Lampard disallowed.[86][87][88]
+In the quarterfinals, Germany defeated Argentina 4–0,[89] and Miroslav Klose tied German Gerd Müller's record of 14 World Cup goals.[90]
+In the semi-final, Germany lost 1–0 to Spain.[91] Germany defeated Uruguay 3–2 to take third place (their second third place after 2006).[92]
+German Thomas Müller won the Golden Boot and the Best Young Player Award.[93][94]
+
+In Euro 2012, Germany was placed in group B along with Portugal, Netherlands, and Denmark. Germany won all three group matches. Germany defeated Greece in the quarter-final and set a record of 15 consecutive wins in all competitive matches.[95] In the semi-finals, Germany lost to Italy, 1–2.
+
+Germany finished first in their qualification group for the 2014 World Cup. The draw for the 2014 World Cup finals placed Germany in Group G,[96] with Portugal, Ghana, and United States. They first faced Portugal in a match billed by some as the "team of all the talents against the team of The Talent (Cristiano Ronaldo)", routing the Portuguese 4–0 thanks to a hat-trick by Thomas Müller.[97][98] In their match with Ghana, they led the game with Götze's second half goal, but then conceded two consecutive goals, then at the 71st minute Klose scored a goal to help Germany to draw 2–2 with Ghana. With that goal, Klose also nudged home his 15th World Cup goal to join former Brazil striker Ronaldo at the pinnacle of World Cup Finals scorers. They then went on to defeat the United States team 1–0, securing them a spot in the round of sixteen against Algeria.
+
+The round of sixteen knockout match against Algeria remained goalless after regulation time, resulting in extra time. In the 92nd minute, André Schürrle scored a goal from a Thomas Müller pass. Mesut Özil scored Germany's second goal in the 120th minute. Algeria managed to score one goal in injury time and the match ended 2–1. Germany secured a place in the quarter-final, where they would face France.
+
+In the quarter-final match against France, Mats Hummels scored in the 13th minute. Germany won the game 1–0 to advance to a record fourth consecutive semi-finals.[99]
+
+The semi-final win (7–1) against Brazil was a major accomplishment. Germany scored four goals in just less than seven minutes and were 5–0 up against Brazil by the 30th minute with goals from Thomas Müller, Miroslav Klose, Sami Khedira and two from Toni Kroos. Klose's goal in the 23rd minute, his 16th World Cup goal, gave him sole possession of the record for most goals scored during World Cup Finals, dethroning former Brazil national Ronaldo.
+
+In the second half of the game, substitute André Schürrle scored twice for Germany to lead 7–0, the highest score against Brazil in a single game. Germany did, however, concede a late goal to Brazil's Oscar. It was Brazil's worst ever World Cup defeat,[100] whilst Germany broke multiple World Cup records with the win, including the record broken by Klose, the first team to reach four consecutive World Cup semi-finals, the first team to score seven goals in a World Cup Finals knockout phase game, the fastest five consecutive goals in World Cup history (four of which in just 400 seconds), the first team to score five goals in the first half in a World Cup semi-final as well as being the topic of the most tweets ever on Twitter about a certain subject when the previous social media record was smashed after Germany scored their fourth goal. Also, Germany's seven goals took their total tally in World Cup history to 223, surpassing Brazil's 221 goals to first place overall.[101]
+
+The World Cup Final was held at the Maracana in Rio de Janeiro on 13 July, and billed as the world's best player (Lionel Messi) versus the world's best team (Germany).[102][103] Mario Götze's 113th-minute goal helped Germany beat Argentina 1–0, becoming the first-ever European team to win a FIFA World Cup in the Americas and the second European team to win the title outside Europe.[104][105]
+
+After several players retired from the team following the 2014 World Cup win, including Philipp Lahm, Per Mertesacker and Miroslav Klose, the team had a disappointing start in the UEFA Euro 2016 qualifiers. They defeated Scotland 2–1 at home, then suffered a 2–0 loss at Poland (the first in their history), a 1–1 draw against the Republic of Ireland, and a 4–0 win over Gibraltar. The year ended with an away 0–1 friendly win against Spain, the reigning European champions of 2008 and 2012.
+
+Troubles during qualifying for the 2016 European Championship continued, drawing at home, as well as losing away, to Ireland; the team also only narrowly defeated Scotland on two occasions, but handily won the return against Poland and both games against Gibraltar (who competed for the first time). Eventually, however, topping their group and qualifying for the tournament through a 2–1 victory against Georgia on 11 October 2015 (having won the first match against them).
+
+On 13 November 2015, the team was playing a friendly match against France in Paris when a series of terrorist attacks took place in the city, some in the direct vicinity of the Stade de France, where the game was held.[106] For security reasons, the team needed to spend the night inside the stadium, accompanied by the French squad who stayed behind in an act of comradery.[107] Four days later, on 17 November 2015, the German team was scheduled to face the Netherlands at Hanover's HDI-Arena, also in a friendly. After initial security reservations, the DFB decided to play the match on 15 November.[108] However, after reports about a concrete threat to the stadium, the match was cancelled ninety minutes before kickoff.[109]
+
+Germany began their preparations for Euro 2016 in March with friendlies against England and Italy. They gave up a 2–0 lead to England, and ended up losing 2–3. They bounced back in their match with Italy, however, winning by a score of 4–1. It was their first win against the Italians in 21 years.[110]
+
+Germany began their campaign for a fourth European title with a 2–0 win against Ukraine on 12 June. Against Poland, Germany was held to a 0–0 draw but concluded Group C with a 1–0 win against Northern Ireland. In the Round of 16, Germany faced Slovakia and earned a comfortable 3–0 win. Germany then faced off against rivals Italy in the quarter-finals. Mesut Özil opened the scoring in the 65th minute for Germany, before Leonardo Bonucci drew even after converting a penalty in the 78th minute. The score remained 1–1 after extra time and Germany beat Italy 6–5 in a penalty shootout. It was the first time Germany had overcome Italy in a major tournament.[111][112] In the semi-finals Germany played the host nation France. Germany's hopes of securing a fourth European championship were put on hold however as France ended Germany's run by eliminating them by a score of 0–2. It was France's first competitive win against Germany in 58 years.[113]
+
+On 2 July 2017, Germany won the 2017 FIFA Confederations Cup after a 1–0 win against Chile in the final at the Krestovsky Stadium in Saint Petersburg, it was their first FIFA Confederations Cup title.[114]
+
+Despite winning all their qualifying matches and the Confederations Cup the previous year, Germany started their 2018 World Cup campaign with a defeat to Mexico. This was their first loss in an opening match since the 1982 World Cup.[115] Germany defeated Sweden 2–1 in their second game via an injury-time winner from Toni Kroos, but was subsequently eliminated following a 2–0 loss to South Korea, their first exit in the first round since 1938 and first ever in group stage since the format had been reintroduced in 1950.[116][117]
+
+Following the World Cup, Germany's struggles continued into the UEFA Nations League. After a 0–0 draw at home against France, they lost 3–0 against the Netherlands[118] and 1–2 in the rematch against France three days later; the latter result being their fourth loss in six competitive matches.[119] These results mean that Germany cannot advance to the 2019 UEFA Nations League Finals and faced the prospect of possible relegation to League B in the next Nations League.[119]
+
+After the Netherlands' win against France, the relegation to League B was originally confirmed, but due to the overhaul of the format for the 2020–21 UEFA Nations League, Germany was spared from relegation to League B.[120]
+
+The national team's home kit has always been a white shirt, black shorts, and white socks. The colours are derived from the 19th-century flag of the North German State of Prussia.[121] Since 1988, many of the home kit's designs incorporate details patterned after the modern German flag. For the 2014 World Cup, the German team used white shorts rather than the traditional black due to FIFA's kit clashing rule for the tournament.[122] The away shirt colour has changed several times. Historically, green shirt with white shorts is the most often used alternative colour combination, derived from the DFB colours – though it is often erroneously reported that the choice is in recognition of the fact that Ireland, whose home shirts are green, were the first nation to play Germany in a friendly game after World War II. However, the first team to play Germany after WWII, as stated above, was actually Switzerland.[123] Other colours such as red, grey and black have also been used.
+
+A change from black to red came in 2005 on the request of Jürgen Klinsmann,[124] but Germany played every game at the 2006 World Cup in its home white colours. In 2010, the away colours then changed back to a black shirt and white shorts, but at the tournament, the team dressed up in the black shorts from the home kit. The German team next resumed the use of a green shirt on its away kit, but then changed again to red-and-black striped shirts with white stripes and letters and black shorts.
+
+Adidas AG is the longstanding kit provider to the national team, a sponsorship that began in 1954 and is contracted to continue until at least 2022.[125] In the 70s, Germany wore Erima kits (a German brand, formerly a subsidiary of Adidas).[126][127]
+
+Germany plays its home matches among various stadiums, in rotation, around the country. They have played home matches in 43 different cities so far, including venues that were German at the time of the match, such as Vienna, Austria, which staged three games between 1938 and 1942.
+
+National team matches have been held most often (46 times) in the stadiums of Berlin, which was the venue of Germany's first home match (in 1908 against England). Other common host cities include Hamburg (34 matches), Stuttgart (32), Hanover (28) and Dortmund. Another notable location is Munich, which has hosted numerous notable matches throughout the history of German football, including the 1974 FIFA World Cup Final, which Germany won against the Netherlands.
+
+Germany's qualifying and friendly matches are televised by privately owned RTL; Nations League by public broadcasters ARD and ZDF. World Cup & European Championships matches featuring the German national team are among the most-watched events in the history of television in Germany.
+
+Recent results and scheduled matches according to the DFB,[133][134] UEFA[135] and FIFA[136] websites.
+
+Win
+  Draw
+  Loss
+  Fixtures
+
+Germany has won the World Cup four times, behind only Brazil (five titles).[137] It has finished as runners-up four times.[137] In terms of semi-final appearances, Germany leads with 13, two more than Brazil's 11, which had participated in two more tournaments.[137] From 1954 to 2014 (16 tournament editions), Germany always reached at least the stage of the last eight teams, before being eliminated in the group stage in 2018.[137] Germany has also qualified for every one of the 18 World Cups for which it has entered – it did not enter the inaugural competition in Uruguay of 1930 for economic reasons, and could not qualify for or compete in the post-war 1950 World Cup as the DFB was reinstated as a FIFA member only two months after this tournament. Germany also has the distinction of having an Elo football rating of 2196 following their victory in the 2014 World Cup, which was higher than any previous champion.[138]
+
+Germany has also won the European Championship three times (Spain and France are the only other multiple-time winners with three and two titles respectively), and finished as runners-up three times as well.[139] The Germans have qualified for every European Championship tournament except for the very first European Championship they entered in 1968.[139] For that tournament, Germany was in the only group of three teams and thus only played four qualifying games. The deciding game was a scoreless draw in Albania which gave Yugoslavia the edge, having won in their neighbour country. The team finished out of top eight only in two occasions, the tournaments of 2000[140] and 2004.[141] In the other ten editions Germany participated in they reached nine times at least the semi-finals, an unparalleled record in Europe.
+
+See also East Germany and Saarland for the results of these separate German teams, and Austria for the team that was merged into the German team from 1938 to 1945.
+
+Champions       Runners-up       Third place       Fourth place
+
+
+
+
+
+FIFA World Cup
+
+UEFA European Championship
+
+FIFA Confederations Cup
+
+
+
+
+
+Source:[145]
+
+The following players were selected for the Euro 2020 qualifying matches against Belarus and Northern Ireland on 16 and 19 November 2019.[146]
+Caps and goals correct as of: 19 November 2019, after the match against Northern Ireland.[147]
+
+The following players have also been called up to the Germany squad within the last 12 months and are still available for selection.
+
+INJ Player withdrew from the squad due to an injury.
+
+Below is a list of the 10 players with the most caps for Germany, as of 22 March 2017[update].[11] (Bold denotes players still available for selection). Players who had played for the separate East German Team (in the scope of this list: Streich 102) do not appear in this list.
+
+Below is a list of the top 10 goalscorers for Germany, as of 27 June 2018[update].[12] (bold denotes players still available for selection). Former East Germany player Joachim Streich, who scored 55 goals, is not included in this Wikipedia list, though he is included in DFB records.

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+
+
+The Argentina national football team (Spanish: Selección de fútbol de Argentina) represents Argentina in men's international football and is administered by the Argentine Football Association, the governing body for football in Argentina. Argentina's home stadium is Estadio Monumental Antonio Vespucio Liberti in Buenos Aires.
+
+La Selección (national team), also known as the La Albiceleste, has appeared in five World Cup finals, including the first final in 1930, which they lost 4–2 to Uruguay. Argentina won in their next final appearance in 1978, beating the Netherlands at extra time, 3–1. Argentina won again in 1986, through a 3–2 victory over West Germany, and a tournament campaign led by Diego Maradona. They made the World Cup finals once more in 1990, and lost 1–0 to West Germany following a controversial penalty call in the 87th minute. Argentina, led by Lionel Messi, made their fifth appearance in a World Cup final in 2014, again losing to Germany, 1–0 during extra-time. Argentina's World Cup winning managers are César Luis Menotti in 1978 and Carlos Bilardo in 1986.
+
+Argentina has also been very successful in the Copa América, winning it 14 times, second only to Uruguay. The team also won the 1992 FIFA Confederations Cup and the 1993 Artemio Franchi Trophy. Argentina is known for having rivalries with Brazil, Uruguay, England, and Germany due to particular occurrences with one another throughout football history.[5][6]
+
+The first match ever recorded for Argentina was against Uruguay.[note 1] The game was held in Montevideo on 16 May 1901 and Argentina won 3–2. During the first years of its existence, the Argentina national team only played friendly matches against other South American teams. The reasons for this varied, including long travel times between countries and the interruption of World War I.[9]
+
+La Selección (national team), also known as the Albicelestes (sky blue and whites), has appeared in five World Cup finals, including the first final in 1930, which they lost, 4–2, to Uruguay. Argentina won in their next final in 1978, beating the Netherlands, 3–1. Argentina, led by Diego Maradona won again in 1986, a 3–2 victory over West Germany.
+Argentina last reached the World Cup final in 2014, where it lost 1–0 to Germany national football team.
+Previous to this their last World Cup final was in 1990, which it also lost, 1–0, to West Germany by a much disputed penalty. Argentina's World Cup winning managers are César Luis Menotti in 1978, and Carlos Bilardo in 1986.
+
+Argentina has been very successful in the Copa América, winning it 14 times. The team also won the FIFA Confederations Cup and the Kirin Cup, both in 1992, and the 1993 Artemio Franchi Trophy. An Argentina team (with only three players of over 23 years of age included in the squad) won the Olympics football tournaments in Athens 2004 and Beijing 2008.[10]
+
+Argentina also won seven of the 18 football competitions at the Pan American Games, winning in 1951, 1955, 1959, 1971, 1995, 2003 and 2019 .
+
+In March 2007, Argentina reached the top of the FIFA World Rankings for the first time.[11]
+
+The River Plate stadium, Monumental Antonio Vespucio Liberti, is a national stadium of Argentina national team that plays most qualifying and friendlies at that stadium.
+
+They play their matches outside the stadium at Córdoba, Rosario, Mendoza, La Plata, San Juan and Salta.
+
+The kit first worn by Argentina was a white shirt, at the official debut of the national side against Uruguay in 1902.[12] In August 1908, Argentina debuted the light blue vertical stripe on white jersey.[13] That kit would become the official kit. The away kits usually have been in dark blue shades, varying the colors of shorts and socks.
+
+Argentina has sported other kits until the blue strip on white kit was made official. On 3 June 1919 in Rio de Janeiro playing the "Roberto Chery Cup" against Brazil, Argentina wore a light blue kit, similar to Uruguay.[14] The trophy was established by Brazilian Football Confederation for the benefit of Roberto Chery's relatives. Chery was Uruguay's substitute goalkeeper and died during the 1919 South American Championship after collapsing in a game against Chile.[15]
+
+At the 1958 World Cup, Argentina wore the yellow jersey of Swedish club IFK Malmö in the match against West Germany, as the team arrived in Sweden without an away kit.
+
+A last moment jersey changed at the 1986 World Cup in Mexico is memorable. Then manager Carlos Bilardo asked the team kit supplier Le Coq Sportif for a lighter blue shirt for the quarter-final in three days against England, that could not be provided. A member of coaching staff scour the shops of Mexico City for 38 shirt plain shirts. They were transformed with an improvised version of the AFA emblem embroidered on to the shirts,[16] and silvery American football numbers ironed to the backs.[17] Argentina beat England with Diego Maradona's "goal of the century".[18][19] The shirt style became an emblem of the occasion and a collector's item.[20]
+
+At the 2018 World Cup in Russia, Argentina debuted a black away kit, a first in their history.[21]
+
+The Argentine Football Association ("AFA") logo has been always used as the team emblem. It debuted in the 1958 World Cup held in Sweden, when Argentina added the AFA logo to their jackets, but not to the shirts.[28]
+
+Nevertheless, the AFA emblem was not used on jerseys until 16 November 1976, when Argentina played the Soviet Union at Estadio Monumental. The first emblem was a simplified version of the crest (without the laurel wreath,[29] that was added for the 1982 World Cup).[16]
+
+In 2004, the two stars added above the crest symbolized the national team FIFA World championships of 1978 and 1986.[29]
+
+Below is a result summary of all matches Argentina have played against FIFA recognized teams.[30]
+
+Positive Record
+  Neutral Record
+  Negative Record
+
+Win
+  Draw
+  Loss
+
+The following players were selected for the 2022 FIFA World Cup qualifiers against Ecuador and Bolivia on 27 and 31 March 2020, respectively. A complementary list including Argentine Primera División players will be released at a further date.[33]
+
+On 12 March 2020, the FIFA announced that the matches originally scheduled to take place during the international window of 23–31 March 2020 are postponed to later dates. Details of the postponed matches will be discussed and announced soon.[34]
+Caps and goals correct as of: 18 November 2019, after the match against Uruguay.
+
+The following players have been called up for the team in the last 12 months.
+
+INJ Withdrew due to injury
+PRE Preliminary squad
+RET Retired from the national team
+SUS Suspended
+
+
+
+
+
+The first Argentina national team manager was Ángel Vázquez, appointed in 1924. Guillermo Stábile is the manager with the most matches coaching the team (127).[51] Here is the complete list of managers:[52][53][54][55]
+
+Argentina have a long and fierce rivalry with their South American neighbours.[56]
+
+With a rivalry stemming from the 1966 World Cup and intensified by the Falklands War of 1982, Argentina and England have had numerous confrontations in World Cup tournaments. Among them was the quarter-final match in 1986, where Diego Maradona scored two goals against England. The first was a handball, but was ruled legal by the referee. The second, scored minutes later, saw Maradona passing five England outfield players before scoring, and is often described as one of the greatest goals in football history.
+
+The nations were paired together in the Round of 16 at the 1998 FIFA World Cup, won by Argentina on penalties, and again at the group stage in 2002, England winning 1–0 through a penalty by David Beckham who had been sent off in the tie four years earlier.
+
+Argentina have played Germany in seven FIFA World Cup matches including three FIFA World Cup finals: In 1986 Argentina won 3–2, but in 1990 it was the Germans who were the victors by a 1–0 scoreline.
+
+In 1958 they met for the first time in the group stage, where Argentina suffered a 1–3 loss to defending champions West Germany.[57] In 1966 both again faced each other in the group stage which ended in a scoreless draw.[58] 2006 they met in the quarter-finals; Argentina lost on penalties after a 1–1 draw. They met again at the same stage in 2010, this time ending with a 4–0 victory for Germany. They played each other for the third consecutive World Cup in the Brazil 2014 event's final, where Argentina were defeated in extra time by a score of 1–0.
+
+Argentina have a long-standing rivalry with their neighbors, that came into existence from the early South American Championships, the 1928 Summer Olympics and the first World Cup final, held in 1930.
+
+Argentina and Uruguay hold the record for most international matches played between two countries.[2] The two teams have faced each other 198 times since 1901. The first match between Argentina and Uruguay was also the first official international match to be played outside the United Kingdom.[note 7]
+
+A minor rivalry developed from the 1990s between Argentina and Nigeria, based not on geographical proximity, long-term battles for honours or factors outside football, but due to the frequency of significant matches between them.[59][60][61][62][63][64] This has included five World Cup group games, all won by Argentina by a single goal margin: 2–1 in 1994, 1–0 in 2002, 1–0 in 2010, 3–2 in 2014 and 2–1 in 2018. The fixture is the most common in the competition's history involving an African nation,[65] and has occurred in five of the six tournaments for which Nigeria has qualified. The sides also met in the 1995 King Fahd Cup (the predecessor to the Confederations Cup) as champions of their respective continents, drawing 0–0.
+
+Below full international level, their Olympic teams also faced off in the gold medal match in 1996 (3–2 to Nigeria), and 2008 (1–0 to Argentina). The final of the 2005 FIFA World Youth Championship was also played between them; both Argentina goals in their 2–1 win were scored by Lionel Messi, who would go on to find the net for the senior team in the 2014[66] and 2018[67] World Cup fixtures. On 6 September 2011, Bangabandhu National Stadium hosted an international friendly football match between the full-strength Argentina and Nigeria teams, featuring Lionel Messi, Sergio Agüero, Javier Mascherano and John Obi Mikel among the other star players of both nations. Argentina won 3–1 with goals from then-Real Madrid teammates Gonzalo Higuaín and Ángel Di María, and an own goal from Nigeria's Elderson Echiéjilé with Chinedu Obasi scoring Nigeria's lone goal.
+
+The sense of rivalry is more keenly felt on the Nigerian side, as Argentina have won almost all of their encounters and have more important traditional opponents to concentrate on, in contrast to the West Africans who remain keen to finally overcome a more illustrious foe.[60]
+
+FIFA World Cup
+
+South American Championship / Copa América
+
+FIFA Confederations Cup
+
+Intercontinental Cup of Nations
+
+Panamerican Championship
+
+Summer Olympics
+
+Pan American Games

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+Fashion is a popular aesthetic expression at a particular time and place and in a specific context, especially in clothing, footwear, lifestyle, accessories, makeup, hairstyle, and body proportions.[1] Whereas a trend often connotes a peculiar aesthetic expression and often lasting shorter than a season, fashion is a distinctive and industry-supported expression traditionally tied to the fashion season and collections.[2] Style is an expression that lasts over many seasons and is often connected to cultural movements and social markers, symbols, class, and culture (ex. Baroque, Rococo, etc.). According to sociologist Pierre Bourdieu, fashion connotes "the latest fashion, the latest difference."[3]
+
+Even though they are often used together, the term fashion differs from clothes and costumes, where the first describes the material and technical garment, whereas the second has been relegated to special senses like fancy-dress or masquerade wear. Fashion instead describes the social and temporal system that "activates" dress as a social signifier in a certain time and context. Philosopher Giorgio Agamben connects fashion to the current intensity of the qualitative moment, to the temporal aspect the Greek called kairos, whereas clothes belong to the quantitative, to what the Greek called Chronos.[4]
+
+Exclusive brands aspire for the label haute couture, but the term is technically limited to members of the Chambre Syndicale de la Haute Couture in Paris.[2] It is more aspirational and inspired by art, culture and movement. It is extremely exclusive in nature.
+
+With increasing mass-production of consumer commodities at Lower prices, and with global reach, sustainability has become an urgent issue amongst politicians, brands, and consumers.[5]
+
+Early Western travelers, traveling to India, Persia, Turkey, or China, would frequently remark on the absence of change in fashion in those countries. The Japanese shōgun's secretary bragged (not completely accurately) to a Spanish visitor in 1609 that Japanese clothing had not changed in over a thousand years.[6] However, there is considerable evidence in Ming China of rapidly changing fashions in Chinese clothing.[7] Changes in costume often took place at times of economic or social change, as occurred in ancient Rome and the medieval Caliphate, followed by a long period without significant changes. In 8th-century Moorish Spain, the musician Ziryab introduced to Córdoba[8][unreliable source][9] sophisticated clothing-styles based on seasonal and daily fashions from his native Baghdad, modified by his inspiration. Similar changes in fashion occurred in the 11th century in the Middle East following the arrival of the Turks, who introduced clothing styles from Central Asia and the Far East.[10]
+
+Additionally, there is a long history of fashion in West Africa.[11] The Cloth was used as a form of currency in trade with the Portuguese and Dutch as early as the 16th Century.[11] Locally produced cloth and cheaper European imports were assembled into new styles to accommodate the growing elite class of West Africans and resident gold and slave traders.[11] There was an exceptionally strong tradition of cloth-weaving in Oyo and the areas inhabited by the Igbo people.[11]
+
+The beginning in Europe of continual and increasingly rapid change in clothing styles can be fairly reliably dated. Historians, including James Laver and Fernand Braudel, date the start of Western fashion in clothing to the middle of the 14th century,[12][13] though they tend to rely heavily on contemporary imagery[14] and illuminated manuscripts were not common before the fourteenth century.[15] The most dramatic early change in fashion was a sudden drastic shortening and tightening of the male over-garment from calf-length to barely covering the buttocks,[16] sometimes accompanied with stuffing in the chest to make it look bigger. This created the distinctive Western outline of a tailored top worn over leggings or trousers.
+
+The pace of change accelerated considerably in the following century, and women's and men's fashion, especially in the dressing and adorning of the hair, became equally complex. Art historians are, therefore, able to use fashion with confidence and precision to date images, often to within five years, particularly in the case of images from the 15th century. Initially, changes in fashion led to a fragmentation across the upper classes of Europe of what had previously been a very similar style of dressing and the subsequent development of distinctive national styles. These national styles remained very different until a counter-movement in the 17th to 18th centuries imposed similar styles once again, mostly originating from Ancien Régime France.[17] Though the rich usually led fashion, the increasing affluence of early modern Europe led to the bourgeoisie and even peasants following trends at a distance, but still uncomfortably close for the elites – a factor that Fernand Braudel regards as one of the main motors of changing fashion.[18]
+
+In the 16th century, national differences were at their most pronounced. Ten 16th century portraits of German or Italian gentlemen may show ten entirely different hats. Albrecht Dürer illustrated the differences in his actual (or composite) contrast of Nuremberg and Venetian fashions at the close of the 15th century (illustration, right). The "Spanish style" of the late 16th century began the move back to synchronicity among upper-class Europeans, and after a struggle in the mid-17th century, French styles decisively took over leadership, a process completed in the 18th century.[20]
+
+Though different textile colors and patterns changed from year to year,[21] the cut of a gentleman's coat and the length of his waistcoat, or the pattern to which a lady's dress was cut, changed more slowly. Men's fashions were primarily derived from military models, and changes in a European male silhouette were galvanized in theaters of European war where gentleman officers had opportunities to make notes of different styles such as the "Steinkirk" cravat or necktie.
+
+Though there had been distribution of dressed dolls from France since the 16th century and Abraham Bosse had produced engravings of fashion in the 1620s, the pace of change picked up in the 1780s with increased publication of French engravings illustrating the latest Paris styles. By 1800, all Western Europeans were dressing alike (or thought they were); local variation became first a sign of provincial culture and later a badge of the conservative peasant.[22]
+
+Although tailors and dressmakers were no doubt responsible for many innovations, and the textile industry indeed led many trends, the history of fashion design is generally understood to date from 1858 when the English-born Charles Frederick Worth opened the first authentic haute couture house in Paris. The Haute house was the name established by the government for the fashion houses that met the standards of the industry. These fashion houses have to adhere to standards such as keeping at least twenty employees engaged in making the clothes, showing two collections per year at fashion shows, and presenting a certain number of patterns to costumers.[23] Since then, the idea of the fashion designer as a celebrity in his or her own right has become increasingly dominant.[24]
+
+Although aspects of fashion can be feminine or masculine, some trends are androgynous.[25] The idea of unisex dressing originated in the 1960s when designers such as Pierre Cardin and Rudi Gernreich created garments, such as stretch jersey tunics or leggings, meant to be worn by both males and females. The impact of unisex expands more broadly to encompass various themes in fashion, including androgyny, mass-market retail, and conceptual clothing.[26] The fashion trends of the 1970s, such as sheepskin jackets, flight jackets, duffel coats, and unstructured clothing, influenced men to attend social gatherings without a tuxedo jacket and to accessorize in new ways. Some men's styles blended the sensuality and expressiveness despite the conservative trend, the growing gay-rights movement and an emphasis on youth allowed for a new freedom to experiment with style, fabrics such as wool crepe, which had previously been associated with women's attire was used by designers when creating male clothing.[27]
+
+The four major current fashion capitals are acknowledged to be Paris, Milan, New York City, and London, which are all headquarters to the most significant fashion companies and are renowned for their major influence on global fashion. Fashion weeks are held in these cities, where designers exhibit their new clothing collections to audiences. A succession of major designers such as Coco Chanel and Yves Saint-Laurent have kept Paris as the center most watched by the rest of the world, although haute couture is now subsidized by the sale of ready-to-wear collections and perfume using the same branding.
+
+Modern Westerners have a vast number of choices available in the selection of their clothes. What a person chooses to wear can reflect his or her personality or interests. When people who have high cultural status start to wear new or different clothes, a fashion trend may start. People who like or respect these people become influenced by their style and begin wearing similarly styled clothes. Fashions may vary considerably within a society according to age, social class, generation, occupation, and geography and may also vary over time. If an older person dresses according to the fashion young people use, he or she may look ridiculous in the eyes of both young and older people. The terms fashionista and fashion victim refer to someone who slavishly follows current fashions.
+
+One can regard the system of sporting various fashions as a fashion language incorporating various fashion statements using a grammar of fashion. (Compare some of the work of Roland Barthes.)
+
+In recent years, Asian fashion has become increasingly significant in local and global markets. Countries such as China, Japan, India, and Pakistan have traditionally had large textile industries, which have often been drawn upon by Western designers, but now Asian clothing styles are also gaining influence based on their ideas.[28]
+
+The notion of the global fashion industry is a product of the modern age.[29] Before the mid-19th century, most clothing was custom-made. It was handmade for individuals, either as home production or on order from dressmakers and tailors. By the beginning of the 20th century—with the rise of new technologies such as the sewing machine, the rise of global capitalism and the development of the factory system of production, and the proliferation of retail outlets such as department stores—clothing had increasingly come to be mass-produced in standard sizes and sold at fixed prices.
+
+Although the fashion industry developed first in Europe and America, as of 2017[update], it is an international and highly globalized industry, with clothing often designed in one country, manufactured in another, and sold worldwide. For example, an American fashion company might source fabric in China and have the clothes manufactured in Vietnam, finished in Italy, and shipped to a warehouse in the United States for distribution to retail outlets internationally. The fashion industry has long been one of the largest employers in the United States,[29] and it remains so in the 21st century. However, U.S. employment declined considerably as production increasingly moved overseas, especially to China. Because data on the fashion industry typically are reported for national economies and expressed in terms of the industry's many separate sectors, aggregate figures for the world production of textiles and clothing are difficult to obtain. However, by any measure, the clothing industry accounts for a significant share of world economic output.[30]
+The fashion industry consists of four levels:
+
+These levels consist of many separate but interdependent sectors. These sectors are Textile Design and Production, Fashion Design and Manufacturing, Fashion Retailing, Marketing and Merchandising, Fashion Shows, and Media and Marketing. Each sector is devoted to the goal of satisfying consumer demand for apparel under conditions that enable participants in the industry to operate at a profit.[29][29][29][29]
+
+Fashion trends influenced by several factors, including cinema, celebrities, climate, creative explorations, political, economic, social, and technological. Examining these factors is called a PEST analysis. Fashion forecasters can use this information to help determine the growth or decline of a particular trend.
+
+Politics has played a central role in the development of fashion. For example, First Lady Jacqueline Kennedy was a fashionable icon of the early 1960s who led the formal dressing trend. By wearing a Chanel suit, a structural Givenchy shift dress, or a soft color Cassini coat with large buttons, it created her elegant look and led a delicate trend.[31]
+
+Furthermore, the political revolution also made much impact on the fashion trend. For example, during the 1960s, the economy had become wealthier, the divorce rate was increasing, and the government approved the birth control pill. This revolution inspired the younger generation to rebellion. In 1964, the leg-baring mini-skirt became a significant fashion trend of the 1960s. Given that fashion designers began to experiment with the shapes of garment, loose sleeveless, micro-minis, flared skirts, and trumpet sleeves. In this case, the mini-skirt trend became an icon of the 1960s.
+
+Moreover, the political movement built an impressive relationship with fashion trends. For instance, during the Vietnam war, the youth of America made a movement that affected the whole country. In the 1960s, the fashion trend was full of fluorescent colors, prints patterns, bell-bottom jeans, fringed vests, and skirt became a protest outfit of the 1960s. This trend was called Hippie, and it is still affecting the current fashion trend.[32]
+
+Technology plays a significant role in most aspects of today's society. Technological influences are growing more apparent in the fashion industry. Advances and new developments are shaping and creating current and future trends.
+
+Developments such as wearable technology have become an essential trend in fashion. They will continue with advances such as clothing constructed with solar panels that charge devices and smart fabrics that enhance wearer comfort by changing color or texture based on environmental changes.[33]
+
+The fashion industry is seeing how 3D printing technology has influenced designers such as Iris Van Herpen and Kimberly Ovitz. These designers have been heavily experimenting and developing 3D printed couture pieces. As the technology grows, the 3D printers will become more accessible to designers and eventually, consumers, which could potentially shape the fashion industry entirely.
+
+Internet technology such as online retailers and social media platforms have given way for trends to be identified, marketed, and sold immediately.[34] Styles and trends are easily conveyed online to attract trendsetters. Posts on Instagram or Facebook can quickly increase awareness about new trends in fashion, which subsequently may create high demand for specific items or brands,[35] new "buy now button" technology can link these styles with direct sales.
+
+Machine vision technology has been developed to track how fashions spread through society. The industry can now see the direct correlation on how fashion shows influence street-chic outfits. The effects can now be quantified and provide valuable feedback to fashion houses, designers, and consumers regarding trends.[36]
+
+Military technology has played an essential role in the fashion industry. The camouflage pattern in clothing was developed to help military personnel be less visible to enemy forces. A trend emerged in the 1960s, and camouflage fabric was introduced to streetwear. The camouflage fabric trend disappeared and resurfaced several times since then. Camouflage started to appear in high fashion by the 1990s.[37] Designers such as Valentino, Dior, and Dolce & Gabbana combined camouflage into their runway and ready-to-wear collections.
+
+Fashion relates to social and cultural context of an environment. According to Matika,[38] "Elements of popular culture become fused when a person's trend is associated with a preference for a genre of music…like music, news or literature, fashion has been fused into everyday lives." Fashion is not only seen as pure aesthetic values; fashion is also a medium for performers to create an overall atmosphere and express their opinions altogether through music video. The latest music video ‘Formation’ by Beyoncé, according to Carlos,[39] "The pop star pays homage to her Creole root.... tracing the roots of the Louisiana cultural nerve center from the post-abolition era to present day, Beyoncé catalogs the evolution of the city's vibrant style and its tumultuous history all at once. Atop a New Orleans police car in a red-and-white Gucci high-collar dress and combat boots, she sits among the ruins of Hurricane Katrina, immediately implanting herself in the biggest national debate on police brutality and race relations in modern day."
+
+Runway show is a reflection of fashion trend and a designer's thought. For designer like Vivienne Westwood, runway shows are a platform for her voice on politics and current events. For her AW15 menswear show, according to Water,[40] "where models with severely bruised faces channeled eco-warriors on a mission to save the planet." Another recent example is a staged feminist protest march for Chanel's SS15 show, rioting models chanting words of empowerment with signs like "Feminist but feminine" and "Ladies first." According to Water,[40] "The show tapped into Chanel's long history of championing female independence: founder Coco Chanel was a trailblazer for liberating the female body in the post-WWI era, introducing silhouettes that countered the restrictive corsets then in favour."
+
+With increasing environmental awareness, the economic imperative to "Spend now, think later" is getting increasingly scrutinized.[41] Today's consumer tends to be more mindful about consumption, looking for just enough and better, more durable options. People have also become more conscious of the impact their everyday consumption has on the environment and society, and these initiatives are often described as a move towards sustainable fashion, yet critics argue a circular economy based on growth is an oxymoron, or an increasing spiral of consumption, rather than a utopian cradle-to-cradle circular solution.
+
+In today's linear economical system, manufacturers extract resources from the earth to make products that will soon be discarded in landfills, on the other hand, under the circular model, the production of goods operates like systems in nature, where the waste and demise of a substance becomes the food and source of growth for something new. Companies such as MUD Jeans, which is based in the Netherlands employs a leasing scheme for jeans. This Dutch company "represents a new consuming philosophy that is about using instead of owning," according to MUD's website. The concept also protects the company from volatile cotton prices. Consumers pay €7.50 a month for a pair of jeans; after a year, they can return the jeans to Mud, trade them for a new pair and start another year-long lease, or keep them. MUD is responsible for any repairs during the lease period.[41] Another ethical fashion company, Patagonia set up the first multi-seller branded store on EBay in order to facilitate secondhand sales; consumers who take the Common Threads pledge can sell in this store and have their gear listed on Patagonia.com's "Used Gear" section.[41]
+
+Consumption as a share of gross domestic product in China has fallen for six decades, from 76 percent in 1952 to 28 percent in 2011. China plans to reduce tariffs on a number of consumer goods and expand its 72-hour transit visa plan to more cities in an effort to stimulate domestic consumption.[42]
+
+The announcement of import tax reductions follows changes in June 2015, when the government cut the tariffs on clothing, cosmetics and various other goods by half. Among the changes — easier tax refunds for overseas shoppers and accelerated openings of more duty-free shops in cities covered by the 72-hour visa scheme. The 72-hour visa was introduced in Beijing and Shanghai in January 2013 and has been extended to 18 Chinese cities.[42]
+
+According to reports at the same time, Chinese consumer spending in other countries such as Japan has slowed even though the yen has dropped.[43] There is clearly a trend in the next 5 years that the domestic fashion market will show an increase.
+
+China is an interesting market for fashion retail as Chinese consumers' motivation to shop for fashion items are unique from Western Audiences.[44] Demographics have limited association with shopping motivation, with occupation, income and education level having no impact; unlike in Western Countries. Chinese high-street shoppers prefer adventure and social shopping, while online shoppers are motivated by idea shopping. Another difference is how gratification and idea shopping influence spending over ¥1k per month on fashion items, and regular spending influenced by value shopping.
+
+Consumers of different groups have varying needs and demands. Factors taken into consideration when thinking of consumers' needs include key demographics.[45]
+To understand consumers' needs and predict fashion trends, fashion companies have to do market research [46] There are two research methods: primary and secondary.[47] Secondary methods are taking other information that has already been collected, for example using a book or an article for research. Primary research is collecting data through surveys, interviews, observation, and/or focus groups. Primary research often focuses on large sample sizes to determine customer's motivations to shop.[44]
+
+Benefits of primary research is specific information about a fashion brand's consumer is explored. Surveys are helpful tools; questions can be open-ended or closed-ended. A negative factor surveys and interviews present is that the answers can be biased, due to wording in the survey or on face-to-face interactions. Focus groups, about 8 to 12 people, can be beneficial because several points can be addressed in depth. However, there are drawbacks to this tactic, too. With such a small sample size, it is hard to know if the greater public would react the same way as the focus group.[47] Observation can really help a company gain insight on what a consumer truly wants. There is less of a bias because consumers are just performing their daily tasks, not necessarily realizing they are being observed. For example, observing the public by taking street style photos of people, the consumer did not get dressed in the morning knowing that would have their photo taken necessarily. They just wear what they would normally wear. Through observation patterns can be seen, helping trend forecasters know what their target market needs and wants.
+
+Knowing the needs of the consumers will increase a fashion companies' sales and profits. Through research and studying the consumers' lives the needs of the customer can be obtained and help fashion brands know what trends the consumers are ready for.
+
+Consumption is driven not only by need, the symbolic meaning for consumers is also a factor. Consumers engaging in symbolic consumption may develop a sense of self over an extended period of time as various objects are collected as part of the process of establishing their identity and, when the symbolic meaning is shared in a social group, to communicate their identity to others. For teenagers consumption plays a role in distinguishing the child self from the adult. Researchers have found that the fashion choices of teenagers are used for self-expression and also to recognize other teens who wear similar clothes. The symbolic association of clothing items can link individuals personality and interests, with music as a prominent factor influencing fashion decisions.[48]
+
+The media plays a significant role when it comes to fashion. For instance, an important part of fashion is fashion journalism. Editorial critique, guidelines, and commentary can be found on television and in magazines, newspapers, fashion websites, social networks, and fashion blogs. In recent years, fashion blogging and YouTube videos have become a major outlet for spreading trends and fashion tips, creating an online culture of sharing one's style on a website or Instagram account. Through these media outlets readers and viewers all over the world can learn about fashion, making it very accessible.[49] In addition to fashion journalism, another media platform that is important in fashion industry is advertisement. Advertisements provide information to audiences and promote the sales of products and services. Fashion industry utilizes advertisements to attract consumers and promote its products to generate sales. Few decades ago when technology was still underdeveloped, advertisements heavily relied on radio, magazines, billboards, and newspapers.[50] These days, there are more various ways in advertisements such as television ads, online-based ads using internet websites, and posts, videos, and live streaming in social media platforms.
+
+At the beginning of the 20th century, fashion magazines began to include photographs of various fashion designs and became even more influential than in the past.[51] In cities throughout the world these magazines were greatly sought after and had a profound effect on public taste in clothing. Talented illustrators drew exquisite fashion plates for the publications which covered the most recent developments in fashion and beauty. Perhaps the most famous of these magazines was La Gazette du Bon Ton, which was founded in 1912 by Lucien Vogel and regularly published until 1925 (with the exception of the war years).[52]
+
+Vogue, founded in the United States in 1892, has been the longest-lasting and most successful of the hundreds of fashion magazines that have come and gone. Increasing affluence after World War II and, most importantly, the advent of cheap color printing in the 1960s, led to a huge boost in its sales and heavy coverage of fashion in mainstream women's magazines, followed by men's magazines in the 1990s. One such example of Vogue's popularity is the younger version, Teen Vogue, which covers clothing and trends that are targeted more toward the "fashionista on a budget". Haute couture designers followed the trend by starting ready-to-wear and perfume lines which are heavily advertised in the magazines and now dwarf their original couture businesses. A recent development within fashion print media is the rise of text-based and critical magazines which aim to prove that fashion is not superficial, by creating a dialogue between fashion academia and the industry. Examples of this trend are: Fashion Theory (1997) and Vestoj (2009). Television coverage began in the 1950s with small fashion features. In the 1960s and 1970s, fashion segments on various entertainment shows became more frequent, and by the 1980s, dedicated fashion shows such as Fashion Television started to appear. FashionTV was the pioneer in this undertaking and has since grown to become the leader in both Fashion Television and new media channels. The Fashion Industry is beginning to promote their styles through Bloggers on social media's. Vogue specified Chiara Ferragni as "blogger of the moment" due to the rises of followers through her Fashion Blog, that became popular.[53]
+
+A few days after the 2010 Fall Fashion Week in New York City came to a close, The New Islander's Fashion Editor, Genevieve Tax, criticized the fashion industry for running on a seasonal schedule of its own, largely at the expense of real-world consumers. "Because designers release their fall collections in the spring and their spring collections in the fall, fashion magazines such as Vogue always and only look forward to the upcoming season, promoting parkas come September while issuing reviews on shorts in January", she writes. "Savvy shoppers, consequently, have been conditioned to be extremely, perhaps impractically, farsighted with their buying."[54]
+
+The fashion industry has been the subject of numerous films and television shows, including the reality show Project Runway and the drama series Ugly Betty. Specific fashion brands have been featured in film, not only as product placement opportunities, but as bespoke items that have subsequently led to trends in fashion.[55]
+
+Videos in general have been very useful in promoting the fashion industry. This is evident not only from television shows directly spotlighting the fashion industry, but also movies, events and music videos which showcase fashion statements as well as promote specific brands through product placements.
+
+There are some fashion advertisements that were accused of racism and led to boycotts from the customers. Globally known, Swedish fashion brand H&M faced this issue with one of its children's wear advertisements in 2018. A black child wearing a hoodie with a slogan written as "coolest monkey in the jungle" right at the center was featured in the ad. When it was released, it immediate became controversial and even led to boycott. A lot of people including celebrities posted on social media about their resentments towards H&M and refusal to work with and buy its products. H&M issued a statement saying "we apologise to anyone this may have offended", which seemed insincere to some.[56] 
+Another fashion advertisement regarding racism is from GAP, American worldwide clothing brand. GAP collaborated with Ellen DeGeneres in 2016 for the advertisement. It features playful, four young girls where tall white girl is leaning with her arm on a shorter black girl's head. When this ad was released, some viewers harshly criticized that it underlies passive racism. A representative from The Root, black culture magazine commented on the ad that it portrays the message that the black people are undervalued and seen like props for white people to look better.[57] There were different points of views on this issue, some saying that people are being too sensitive, and some getting offended. Regardless of various views and thoughts, GAP replaced the ad to different image and apologized to critics.[58]
+
+Many fashion brands have published ads that were too provocative and sexy to attract customers’ attention. British high fashion brand, Jimmy Choo, was blamed for having sexism in its ad which featured a female British mode wearing the brand's boots. In this two-minute ad, men whistle at a model, walking on the street with red, sleeveless mini dress. This ad gained much backlash and criticism by the viewers since sexual harassment and misconduct were a huge issue during this time and even till now. Many people showed their dismay through social media posts, leading Jimmy Choo to pull down the ad from social media platforms.[59]
+
+French luxury fashion brand Yves Saint Laurent also faced this issue with its print ad shown in Paris in 2017. A female model is wearing a fishnet tights with roller-skate stilettos, almost lying down with her legs opened in front of the camera. This advertisement brought harsh comments from the viewers and French advertising organization directors for going against the advertising codes related to "respect for decency, dignity and those prohibiting submission, violence or dependence, as well as the use of stereotypes." They even said that this ad is causing "mental harm to adolescents."[60] Lot of sarcastic comments were made in social media about the ad and the poster was removed from the city.
+
+Fashion public relations involves being in touch with a company's audiences and creating strong relationships with them, reaching out to media and initiating messages that project positive images of the company.[61] Social media plays an important role in modern-day fashion public relations; enabling practitioners to reach a wide range of consumers through various platforms.[62]
+
+Building brand awareness and credibility is a key implication of good public relations. In some cases, great hype is built about new designers' collections before they are released into the market, due to the immense exposure generated by practitioners.[63] Social media, such as blogs, micro blogs, podcasts, photo and video sharing sites have all become increasingly important to fashion public relations.[64] The interactive nature of these platforms allows practitioners to engage and communicate with the public in real time, and tailor their clients' brand or campaign messages to the target audience. With blogging platforms such as Instagram, Tumblr, Wordpress, and other sharing sites, bloggers have emerged as expert fashion commentators, shaping brands and having a great impact on what is ‘on trend’.[65] Women in the fashion public relations industry such as Sweaty Betty PR founder Roxy Jacenko and Oscar de la Renta's PR girl Erika Bearman, have acquired copious followers on their social media sites, by providing a brand identity and a behind the scenes look into the companies they work for.
+
+Social media is changing the way practitioners deliver messages,[16] as they are concerned with the media, and also customer relationship building.[66] PR practitioners must provide effective communication among all platforms, in order to engage the fashion public in an industry socially connected via online shopping.[67] Consumers have the ability to share their purchases on their personal social media pages (such as Facebook, Twitter, Instagram, etc.), and if practitioners deliver the brand message effectively and meet the needs of its public, word-of-mouth publicity will be generated and potentially provide a wide reach for the designer and their products.
+
+Anthropology, the study of culture and human societies, studies fashion by asking why certain styles are deemed socially appropriate and others are not. A certain way is chosen and that becomes the fashion as defined by a certain people as a whole, so if a particular style has a meaning in an already occurring set of beliefs that style will become fashion.[68] According to Ted Polhemus and Lynn Procter, fashion can be described as adornment, of which there are two types: fashion and anti-fashion. Through the capitalization and commoditisation of clothing, accessories, and shoes, etc., what once constituted anti-fashion becomes part of fashion as the lines between fashion and anti-fashion are blurred.[69]
+
+The definition of fashion and anti-fashion is as follows: Anti-fashion is fixed and changes little over time. Anti-fashion is different depending on the cultural or social group one is associated with or where one lives, but within that group or locality the style changes little. Fashion is the exact opposite of anti-fashion. Fashion changes very quickly and is not affiliated with one group or area of the world but is spread out throughout the world wherever people can communicate easily with each other. For example, Queen Elizabeth II's 1953 coronation gown is an example of anti-fashion because it is traditional and does not change over any period whereas a gown from fashion designer Dior's collection of 1953 is fashion because the style will change every season as Dior comes up with a new gown to replace the old one. In the Dior gown the length, cut, fabric, and embroidery of the gown change from season to season. Anti-fashion is concerned with maintaining the status quo while fashion is concerned with social mobility. Time is expressed in terms of continuity in anti-fashion and as change in fashion. Fashion has changing modes of adornment while anti-fashion has fixed modes of adornment. Indigenous and peasant modes of adornment are an example of anti-fashion. Change in fashion is part of the larger system and is structured to be a deliberate change in style.[70]
+
+Today, people in rich countries are linked to people in poor countries through the commoditization and consumption of what is called fashion. People work long hours in one area of the globe to produce things that people in another part of the globe are anxious to consume. An example of this is the chain of production and consumption of Nike shoes, which are produced in Taiwan and then purchased in North America. At the production end, there is nation-building a hard working ideology that leads people to produce and entices people to consume with a vast amount of goods for the offering[clarification needed]. Commodities are no longer just utilitarian but are fashionable, be they running shoes or sweat suits.[71]
+
+The change from anti-fashion to fashion because of the influence of western consumer-driven civilization can be seen in eastern Indonesia. The ikat textiles of the Ngada area of eastern Indonesia are changing because of modernization and development. Traditionally, in the Ngada area there was no idea similar to that of the Western idea of fashion, but anti-fashion in the form of traditional textiles and ways to adorn oneself were widely popular. Textiles in Indonesia have played many roles for the local people. Textiles defined a person's rank and status; certain textiles indicated being part of the ruling class. People expressed their ethnic identity and social hierarchy through textiles. Because some Indonesians bartered ikat textiles for food, the textiles constituted economic goods, and as some textile design motifs had spiritual religious meanings, textiles were also a way to communicate religious messages.[72]
+
+In eastern Indonesia, both the production and use of traditional textiles have been transformed as the production, use and value associated with textiles have changed due to modernization. In the past, women produced the textiles either for home consumption or to trade with others. Today, this has changed as most textiles are not being produced at home. Western goods are considered modern and are valued more than traditional goods, including the sarong, which retain a lingering association with colonialism. Now, sarongs are used only for rituals and ceremonial occasions, whereas western clothes are worn to church or government offices. Civil servants working in urban areas are more likely than peasants to make the distinction between western and traditional clothes. Following Indonesia's independence from the Dutch, people increasingly started buying factory made shirts and sarongs. In textile-producing areas the growing of cotton and production of naturally colored thread became obsolete. Traditional motifs on textiles are no longer considered the property of a certain social class or age group. Wives of government officials are promoting the use of traditional textiles in the form of western garments such as skirts, vests and blouses. This trend is also being followed by the general populace, and whoever can afford to hire a tailor is doing so to stitch traditional ikat textiles into western clothes. Thus, traditional textiles are now fashion goods and are no longer confined to the black, white and brown colour palette but come in array of colours. Traditional textiles are also being used in interior decorations and to make handbags, wallets and other accessories, which are considered fashionable by civil servants and their families. There is also a booming tourist trade in the eastern Indonesian city of Kupang where international as well as domestic tourists are eager to purchase traditionally printed western goods.[73]
+
+The use of traditional textiles for fashion is becoming big business in eastern Indonesia, but these traditional textiles are losing their ethnic identity markers and are being used as an item of fashion.[74]
+
+In the fashion industry, intellectual property is not enforced as it is within the film industry and music industry. Robert Glariston, an intellectual property expert, mentioned in a fashion seminar held in LA[which?] that "Copyright law regarding clothing is a current hot-button issue in the industry. We often have to draw the line between designers being inspired by a design and those outright stealing it in different places."[75] To take inspiration from others' designs contributes to the fashion industry's ability to establish clothing trends. For the past few years, WGSN has been a dominant source of fashion news and forecasts in encouraging fashion brands worldwide to be inspired by one another. Enticing consumers to buy clothing by establishing new trends is, some have argued, a key component of the industry's success. Intellectual property rules that interfere with this process of trend-making would, in this view, be counter-productive. On the other hand, it is often argued that the blatant theft of new ideas, unique designs, and design details by larger companies is what often contributes to the failure of many smaller or independent design companies.
+
+Since fakes are distinguishable by their poorer quality, there is still a demand for luxury goods, and as only a trademark or logo can be copyrighted, many fashion brands make this one of the most visible aspects of the garment or accessory. In handbags, especially, the designer's brand may be woven into the fabric (or the lining fabric) from which the bag is made, making the brand an intrinsic element of the bag.
+
+In 2005, the World Intellectual Property Organization (WIPO) held a conference calling for stricter intellectual property enforcement within the fashion industry to better protect small and medium businesses and promote competitiveness within the textile and clothing industries.[76][77]
+
+There has been great debate about politics' place in fashion and traditionally, the fashion industry has maintained a rather apolitical stance.[78] Considering the U.S.'s political climate in the surrounding months of the 2016 presidential election, during 2017 fashion weeks in London, Milan, New York, Paris and São Paulo amongst others, many designers took the opportunity to take political stances leveraging their platforms and influence to reach the masses.[79][80]
+
+Aiming to "amplify a greater message of unity, inclusion, diversity, and feminism in a fashion space", Mara Hoffman invited the founders of the Women's March on Washington to open her show which featured modern silhouettes of utilitarian wear, described by critics as "Made for a modern warrior" and "Clothing for those who still have work to do".[81] Prabal Gurung debuted his collection of T-shirts featuring slogans such as "The Future is Female", "We Will Not Be Silenced", and "Nevertheless She Persisted", with proceeds going to the ACLU, Planned Parenthood, and Gurung's own charity, "Shikshya Foundation Nepal".[78] Similarly, The Business of Fashion launched the #TiedTogether movement on Social Media, encouraging member of the industry from editors to models, to wear a white bandana advocating for "unity, solidarity, and inclusiveness during fashion week".[82]
+
+Fashion may be used to promote a cause, such as to promote healthy behavior,[83] to raise money for a cancer cure,[84] or to raise money for local charities[85] such as the Juvenile Protective Association[86] or a children's hospice.[87]
+
+One fashion cause is trashion, which is using trash to make clothes, jewelry, and other fashion items in order to promote awareness of pollution. There are a number of modern trashion artists such as Marina DeBris, Ann Wizer,[88] and Nancy Judd.[89]
+
+African-Americans have used fashion through the years, to express themselves and their ideas.[90] It has grown and developed with time. African-American influencers often have been known to start trends though modern day social media, and even in past years they have been able to reach others with their fashion and style.
+
+Celebrities like Rihanna, Lupita Nyong'o, Zendaya, and Michelle Obama have been a few of the many fashion idols in the black female community. For men, Pharrell Williams, Kanye West, and Ice Cube have also helped define modern day fashion for black men. Today's fashion scene is not just clothes, but also hair and makeup. Recent trends have included the embracing of natural hair, traditional clothing worn with modern clothing, or traditional patterns used in modern clothing styles. All of these trends come with the long existing and persevering movement of "Black is Beautiful".
+
+In the mid to end of the 1900s, African American style changed and developed with the times. Around the 1950s is really when the black community was able to create their own distinct styles. The term "Sunday attire" was coined, communities emphasized "Correct" dress, it was especially important when "stepping out" for social occasions with community members, a habit that continues in the early 2000s.[91] Hair-dos and hairstyles also became a fashion statement, for example the "conk" which is hair that is slightly flattened and waved.[91] Afros also emerged and they were often used to symbolize the rejection of white beauty standards at the time.[92] Around the 1970s is when flashy costumes began to appear and black artists really started to define their presences through fashion. Around this time is also when movements started using fashion as one of their outlets.[92]
+
+Black activists and supporters used fashion to express their solidarity and support of this civil rights movement. Supporters adorned symbolic clothing, accessories and hairstyles, usually native to Africa. Politics and fashion were fused together during this time and the use of these symbolic fashion statements sent a message to America and the rest of the world that African Americans were proud of their heritage.[92] They aimed to send an even stronger message that black is beautiful and they were not afraid to embrace their identities.[92] An example would the Kente cloth, it is a brightly colored strip of cloth that is stitched and woven together to create different accessories.[92] This woven cloth of brightly colored strips of fabric became a strong symbolic representation of pride in African identity for African Americans of the 1960s and later. It was developed into what is called a dashiki, a flowing, loose fitting, tunic style shirt. This cloth became one of the most notorious symbols of this revolution.[93]
+
+The Black Panther Party (BPP) was an essential piece of the Black Power movement that allowed members that were involved advocate for the African American race in different subjects like equality and politics. The BPP members wore a very distinctive uniform: a black leather jacket, black pants, light blue shirts, a black beret, an afro, dark sunglasses, and usually a fist in the air.[94] Their image gave off a very militant like feel to it. This notable uniform was established in 1996, but a different uniform was still in place before; just the sunglasses and leather jackets.[94] Each member wore this uniform at events, rallies, and in their day-today life. Very few members changed the essential parts of the outfit, but some added personal touches such as necklaces or other jewelry that was usually were a part of African culture.[93] The Black Panther uniform did succeeded in intimidating enemies and onlookers and clearly sent a message of black pride and power even though the initial intention of this party was to communicate solidarity among the Black Panther Party members.[94]
+
+Since the 1970s, fashion models of color, especially black men and women, have experienced an increase in discrimination in the fashion industry. In the years from 1970 to 1990, black designers and models were very successful, but as the 1990s came to an end, the fashion aesthetic changed and it did not include black models or designers.[95] In today's fashion, black models, influencers, and designers account for one of the smallest percentages of the industry.[95] There are many theories about this lack of diversity, that it can be attributed to the economic differences usually associated with race and class, or it can reflect the differences in arts schooling given to mostly black populated schools, and also blatant racism.
+
+A report from New York Fashion (Spring 2015) week found that while 79.69% of models on the runway were white, only 9.75% of models were black, 7.67% were Asian, and 2.12% were Latina. The lack of diversity also accounts for not only designers but models too, out of four hundred and seventy members of The Council of Fashion Designers of America (CFDA) only twelve of the members are black.[96] From the same study on New York Fashion Week, it was shown that only 2.7% of the 260 designers presented were black men, and an even smaller percentage were black female designers.[96] Even the relationship between independent designers and retailers can show the racial gap, only 1% of designers stocked at department stores being people of color. It was also found that in editorial spreads, over eighty percent of models pictured were white and only nine percent were black models. These numbers have stayed stagnant over the past few years.[96]
+
+Many fashion designers have come under fire over the years for what is known as tokenism. Designer or editors will add one or two members on an underrepresented group to help them appear as inclusive and diverse, and to also help them give the illusion that they have equality.[95] This idea of tokenism helps designers avoid accusations of racism, sexism, body shaming, etc.[95]
+
+There are many examples of cultural appropriation in fashion. In many instances, designers can be found using aspects of culture inappropriately, in most cases taking traditional clothing from middle eastern, African, and Hispanic culture and adding it to their runway fashion.[97] Some examples are in a 2018 Gucci runway show, white models wore Sikh headdresses, causing a lot of backlash. Victoria's secret was also under fire for putting traditional native headdresses on their models during a lingerie runway show. Marc Jacobs sent down models sporting dreadlocks in his spring 2017 New York Fashion Week show, also facing immense criticism.[98]

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+Aluminium (aluminum in American and Canadian English) is a chemical element with the symbol Al and atomic number 13. It is a silvery-white, soft, non-magnetic and ductile metal in the boron group. By mass, aluminium makes up about 8% of the Earth's crust, where it is the third most abundant element (after oxygen and silicon) and also the most abundant metal. Occurrence of aluminium decreases in the Earth's mantle below, however. The chief ore of aluminium is bauxite. Aluminium metal is highly reactive, such that native specimens are rare and limited to extreme reducing environments. Instead, it is found combined in over 270 different minerals.[7]
+
+Aluminium is remarkable for its low density and its ability to resist corrosion through the phenomenon of passivation. Aluminium and its alloys are vital to the aerospace industry[8] and important in transportation and building industries, such as building facades and window frames.[9] The oxides and sulfates are the most useful compounds of aluminium.[8]
+
+Despite its prevalence in the environment, no known form of life uses aluminium salts metabolically, but aluminium is well tolerated by plants and animals.[10] Because of these salts' abundance, the potential for a biological role for them is of continuing interest, and studies continue.
+
+Of aluminium isotopes, only 27Al is stable. This is consistent with aluminium having an odd atomic number.[b] It is the only primordial aluminium isotope, i.e. the only one that has existed on Earth in its current form since the creation of the planet. Nearly all aluminium on Earth is present as this isotope, which makes it a mononuclidic element and means that its standard atomic weight is the same as that of the isotope. The standard atomic weight of aluminium is low in comparison with many other metals,[c] which has consequences for the element's properties (see below). This makes aluminium very useful in nuclear magnetic resonance (NMR), as its single stable isotope has a high NMR sensitivity.[12]
+
+All other isotopes of aluminium are radioactive. The most stable of these is 26Al: while it was present along with stable 27Al in the interstellar medium from which the Solar System formed, having been produced by stellar nucleosynthesis as well, its half-life is only 717,000 years and therefore it could not have survived since the formation of the planet. However, minute traces of 26Al are produced from argon in the atmosphere by spallation caused by cosmic ray protons. The ratio of 26Al to 10Be has been used for radiodating of geological processes over 105 to 106 year time scales, in particular transport, deposition, sediment storage, burial times, and erosion.[13] Most meteorite scientists believe that the energy released by the decay of 26Al was responsible for the melting and differentiation of some asteroids after their formation 4.55 billion years ago.[14]
+
+The remaining isotopes of aluminium, with mass numbers ranging from 22 to 43, all have half-lives well under an hour. Three metastable states are known, all with half-lives under a minute.[11]
+
+An aluminium atom has 13 electrons, arranged in an electron configuration of [Ne] 3s2 3p1,[15] with three electrons beyond a stable noble gas configuration. Accordingly, the combined first three ionization energies of aluminium are far lower than the fourth ionization energy alone.[16] Such an electron configuration is shared with the other well-characterized members of its group, boron, gallium, indium, and thallium; it is also expected for nihonium. Aluminium can relatively easily surrender its three outermost electrons in many chemical reactions (see below). The electronegativity of aluminium is 1.61 (Pauling scale).[17]
+
+A free aluminium atom has a radius of 143 pm.[18] With the three outermost electrons removed, the radius shrinks to 39 pm for a 4-coordinated atom or 53.5 pm for a 6-coordinated atom.[18] At standard temperature and pressure, aluminium atoms (when not affected by atoms of other elements) form a face-centered cubic crystal system bound by metallic bonding provided by atoms' outermost electrons; hence aluminium (at these conditions) is a metal.[19] This crystal system is shared by many other metals, such as lead and copper; the size of a unit cell of aluminium is comparable to that of those other metals.[19] It is however not shared by the other members of its group; boron has ionization energies too high to allow metallization, thallium has a hexagonal close-packed structure, and gallium and indium have unusual structures that are not close-packed like those of aluminium and thallium. Since few electrons are available for metallic bonding, aluminium metal is soft with a low melting point and low electrical resistivity, as is common for post-transition metals.[20]
+
+Aluminium metal has an appearance ranging from silvery white to dull gray, depending on the surface roughness. A fresh film of aluminium serves as a good reflector (approximately 92%) of visible light and an excellent reflector (as much as 98%) of medium and far infrared radiation.
+
+The density of aluminium is 2.70 g/cm3, about 1/3 that of steel, much lower than other commonly encountered metals, making aluminium parts easily identifiable through their lightness.[21] Aluminium's low density compared to most other metals arises from the fact that its nuclei are much lighter, while difference in the unit cell size does not compensate for this difference. The only lighter metals are the metals of groups 1 and 2, which apart from beryllium and magnesium are too reactive for structural use (and beryllium is very toxic).[22] Aluminium is not as strong or stiff as steel, but the low density makes up for this in the aerospace industry and for many other applications where light weight and relatively high strength are crucial.
+
+Pure aluminium is quite soft and lacking in strength. In most applications various aluminium alloys are used instead because of their higher strength and hardness. The yield strength of pure aluminium is 7–11 MPa, while aluminium alloys have yield strengths ranging from 200 MPa to 600 MPa.[23] Aluminium is ductile, with a percent elongation of 50-70%,[24] and malleable allowing it to be easily drawn and extruded. It is also easily machined, and the low melting temperature of 660 °C allows for easy casting.
+
+Aluminium is an excellent thermal and electrical conductor, having 59% the conductivity of copper, both thermal and electrical, while having only 30% of copper's density. Aluminium is capable of superconductivity, with a superconducting critical temperature of 1.2 kelvin and a critical magnetic field of about 100 gauss (10 milliteslas).[25] It is paramagnetic and thus essentially unaffected by static magnetic fields. The high electrical conductivity, however, means that it is strongly affected by changing magnetic field through the induction of eddy currents.
+
+Aluminium combines characteristics of pre- and post-transition metals. Since it has few available electrons for metallic bonding, like its heavier group 13 congeners, it has the characteristic physical properties of a post-transition metal, with longer-than-expected interatomic distances.[20] Furthermore, as Al3+ is a small and highly charged cation, it is strongly polarizing and aluminium compounds tend towards covalency;[26] this behaviour is similar to that of beryllium (Be2+), and the two display an example of a diagonal relationship.[27]
+
+The underlying core under aluminium's valence shell is that of the preceding noble gas, whereas those of its heavier congeners gallium and indium, thallium, and nihonium also include a filled d-subshell and in some cases a filled f-subshell. Hence, the inner electrons of aluminium shield the valence electrons almost completely, unlike those of aluminium's heavier congeners. As such, aluminium is the most electropositive metal in its group. In fact, aluminium's electropositive behavior, high affinity for oxygen, and highly negative standard electrode potential are all more similar to those of scandium, yttrium, lanthanum, and actinium, which like aluminium have three valence electrons outside a noble gas core.[20] Aluminium also bears minor similarities to the metalloid boron in the same group: AlX3 compounds are valence isoelectronic to BX3 compounds (they have the same valence electronic structure), and both behave as Lewis acids and readily form adducts.[28] Additionally, one of the main motifs of boron chemistry is regular icosahedral structures, and aluminium forms an important part of many icosahedral quasicrystal alloys, including the Al–Zn–Mg class.[29]
+
+Aluminium has a high chemical affinity to oxygen, which renders it suitable for use as a reducing agent in the thermite reaction. A fine powder of aluminium metal reacts explosively on contact with liquid oxygen; under normal conditions, however, aluminium forms a thin oxide layer (~ 5 nm at room temperature)[30] that protects the metal from further corrosion by oxygen, water, or dilute acid, a process termed passivation.[26][31] Because of its general resistance to corrosion, aluminium is one of the few metals that retains silvery reflectance in finely powdered form, making it an important component of silver-colored paints.[32] Aluminium is not attacked by oxidizing acids because of its passivation. This allows aluminium to be used to store reagents such as nitric acid, concentrated sulfuric acid, and some organic acids.[10]
+
+In hot concentrated hydrochloric acid, aluminium reacts with water with evolution of hydrogen, and in aqueous sodium hydroxide or potassium hydroxide at room temperature to form aluminates—protective passivation under these conditions is negligible.[33] Aqua regia also dissolves aluminium.[10] Aluminium is corroded by dissolved chlorides, such as common sodium chloride, which is why household plumbing is never made from aluminium.[33] The oxide layer on aluminium is also destroyed by contact with mercury due to amalgamation or with salts of some electropositive metals.[26] As such, the strongest aluminium alloys are less corrosion-resistant due to galvanic reactions with alloyed copper,[23] and aluminium's corrosion resistance is greatly reduced by aqueous salts, particularly in the presence of dissimilar metals.[20]
+
+Aluminium reacts with most nonmetals upon heating, forming compounds such as aluminium nitride (AlN), aluminium sulfide (Al2S3), and the aluminium halides (AlX3). It also forms a wide range of intermetallic compounds involving metals from every group on the periodic table.[26]
+
+The vast majority of compounds, including all aluminium-containing minerals and all commercially significant aluminium compounds, feature aluminium in the oxidation state 3+. The coordination number of such compounds varies, but generally Al3+ is either six- or four-coordinate. Almost all compounds of aluminium(III) are colorless.[26]
+
+In aqueous solution, Al3+ exists as the hexaaqua cation [Al(H2O)6]3+, which has an approximate pKa of 10−5.[12] Such solutions are acidic as this cation can act as a proton donor and progressively hydrolyse until a precipitate of aluminium hydroxide, Al(OH)3, forms. This is useful for clarification of water, as the precipitate nucleates on suspended particles in the water, hence removing them. Increasing the pH even further leads to the hydroxide dissolving again as aluminate, [Al(H2O)2(OH)4]−, is formed.
+
+Aluminium hydroxide forms both salts and aluminates and dissolves in acid and alkali, as well as on fusion with acidic and basic oxides.[26] This behaviour of Al(OH)3 is termed amphoterism, and is characteristic of weakly basic cations that form insoluble hydroxides and whose hydrated species can also donate their protons. One effect of this is that aluminium salts with weak acids are hydrolysed in water to the aquated hydroxide and the corresponding nonmetal hydride: for example, aluminium sulfide yields hydrogen sulfide. However, some salts like aluminium carbonate exist in aqueous solution but are unstable as such; and only incomplete hydrolysis takes place for salts with strong acids, such as the halides, nitrate, and sulfate. For similar reasons, anhydrous aluminium salts cannot be made by heating their "hydrates": hydrated aluminium chloride is in fact not AlCl3·6H2O but [Al(H2O)6]Cl3, and the Al–O bonds are so strong that heating is not sufficient to break them and form Al–Cl bonds instead:[26]
+
+All four trihalides are well known. Unlike the structures of the three heavier trihalides, aluminium fluoride (AlF3) features six-coordinate aluminium, which explains its involatility and insolubility as well as high heat of formation. Each aluminium atom is surrounded by six fluorine atoms in a distorted octahedral arrangement, with each fluorine atom being shared between the corners of two octahedra. Such {AlF6} units also exist in complex fluorides such as cryolite, Na3AlF6.[d] AlF3 melts at 1,290 °C (2,354 °F) and is made by reaction of aluminium oxide with hydrogen fluoride gas at 700 °C (1,292 °F).[35]
+
+With heavier halides, the coordination numbers are lower. The other trihalides are dimeric or polymeric with tetrahedral four-coordinate aluminium centers. Aluminium trichloride (AlCl3) has a layered polymeric structure below its melting point of 192.4 °C (378 °F) but transforms on melting to Al2Cl6 dimers. At higher temperatures those increasingly dissociate into trigonal planar AlCl3 monomers similar to the structure of BCl3. Aluminium tribromide and aluminium triiodide form Al2X6 dimers in all three phases and hence do not show such significant changes of properties upon phase change.[35] These materials are prepared by treating aluminium metal with the halogen. The aluminium trihalides form many addition compounds or complexes; their Lewis acidic nature makes them useful as catalysts for the Friedel–Crafts reactions. Aluminium trichloride has major industrial uses involving this reaction, such as in the manufacture of anthraquinones and styrene; it is also often used as the precursor for many other aluminium compounds and as a reagent for converting nonmetal fluorides into the corresponding chlorides (a transhalogenation reaction).[35]
+
+Aluminium forms one stable oxide with the chemical formula Al2O3, commonly called alumina.[36] It can be found in nature in the mineral corundum, α-alumina;[37] there is also a γ-alumina phase.[12] Its crystalline form, corundum, is very hard (Mohs hardness 9), has a high melting point of 2,045 °C (3,713 °F), has very low volatility, is chemically inert, and a good electrical insulator, it is often used in abrasives (such as toothpaste), as a refractory material, and in ceramics, as well as being the starting material for the electrolytic production of aluminium metal. Sapphire and ruby are impure corundum contaminated with trace amounts of other metals.[12] The two main oxide-hydroxides, AlO(OH), are boehmite and diaspore. There are three main trihydroxides: bayerite, gibbsite, and nordstrandite, which differ in their crystalline structure (polymorphs). Many other intermediate and related structures are also known.[12] Most are produced from ores by a variety of wet processes using acid and base. Heating the hydroxides leads to formation of corundum. These materials are of central importance to the production of aluminium and are themselves extremely useful. Some mixed oxide phases are also very useful, such as spinel (MgAl2O4), Na-β-alumina (NaAl11O17), and tricalcium aluminate (Ca3Al2O6, an important mineral phase in Portland cement).[12]
+
+The only stable chalcogenides under normal conditions are aluminium sulfide (Al2S3), selenide (Al2Se3), and telluride (Al2Te3). All three are prepared by direct reaction of their elements at about 1,000 °C (1,832 °F) and quickly hydrolyse completely in water to yield aluminium hydroxide and the respective hydrogen chalcogenide. As aluminium is a small atom relative to these chalcogens, these have four-coordinate tetrahedral aluminium with various polymorphs having structures related to wurtzite, with two-thirds of the possible metal sites occupied either in an orderly (α) or random (β) fashion; the sulfide also has a γ form related to γ-alumina, and an unusual high-temperature hexagonal form where half the aluminium atoms have tetrahedral four-coordination and the other half have trigonal bipyramidal five-coordination.[38]
+
+Four pnictides – aluminium nitride (AlN), aluminium phosphide (AlP), aluminium arsenide (AlAs), and aluminium antimonide (AlSb) – are known. They are all III-V semiconductors isoelectronic to silicon and germanium, all of which but AlN have the zinc blende structure. All four can be made by high-temperature (and possibly high-pressure) direct reaction of their component elements.[38]
+
+Although the great majority of aluminium compounds feature Al3+ centers, compounds with lower oxidation states are known and are sometimes of significance as precursors to the Al3+ species.
+
+AlF, AlCl, AlBr, and AlI exist in the gaseous phase when the respective trihalide is heated with aluminium, and at cryogenic temperatures. Their instability in the condensed phase is due to their ready disproportionation to aluminium and the respective trihalide: the reverse reaction is favored at high temperature (although even then they are still short-lived), explaining why AlF3 is more volatile when heated in the presence of aluminium metal, as is aluminium metal when heated in the presence of AlCl3.[35] A stable derivative of aluminium monoiodide is the cyclic adduct formed with triethylamine, Al4I4(NEt3)4. Also of theoretical interest but only of fleeting existence are Al2O and Al2S. Al2O is made by heating the normal oxide, Al2O3, with silicon at 1,800 °C (3,272 °F) in a vacuum. Such materials quickly disproportionate to the starting materials.[39]
+
+Very simple Al(II) compounds are invoked or observed in the reactions of Al metal with oxidants. For example, aluminium monoxide, AlO, has been detected in the gas phase after explosion[40] and in stellar absorption spectra.[41] More thoroughly investigated are compounds of the formula R4Al2 which contain an Al–Al bond and where R is a large organic ligand.[42]
+
+A variety of compounds of empirical formula AlR3 and AlR1.5Cl1.5 exist.[43] The aluminium trialkyls and triaryls are reactive, volatile, and colorless liquids or low-melting solids. They catch fire spontaneously in air and react with water, thus necessitating precautions when handling them. They often form dimers, unlike their boron analogues, but this tendency diminishes for branched-chain alkyls (e.g. Pri, Bui, Me3CCH2); for example, triisobutylaluminium exists as an equilibrium mixture of the monomer and dimer.[44][45] These dimers, such as trimethylaluminium (Al2Me6), usually feature tetrahedral Al centers formed by dimerization with some alkyl group bridging between both aluminium atoms. They are hard acids and react readily with ligands, forming adducts. In industry, they are mostly used in alkene insertion reactions, as discovered by Karl Ziegler, most importantly in "growth reactions" that form long-chain unbranched primary alkenes and alcohols, and in the low-pressure polymerization of ethene and propene. There are also some heterocyclic and cluster organoaluminium compounds involving Al–N bonds.[44]
+
+The industrially most important aluminium hydride is lithium aluminium hydride (LiAlH4), which is used in as a reducing agent in organic chemistry. It can be produced from lithium hydride and aluminium trichloride.[46] The simplest hydride, aluminium hydride or alane, is not as important. It is a polymer with the formula (AlH3)n, in contrast to the corresponding boron hydride that is a dimer with the formula (BH3)2.[46]
+
+Aluminium's per-particle abundance in the Solar System is 3.15 ppm (parts per million).[47][e] It is the twelfth most abundant of all elements and third most abundant among the elements that have odd atomic numbers, after hydrogen and nitrogen.[47] The only stable isotope of aluminium, 27Al, is the eighteenth most abundant nucleus in the Universe. It is created almost entirely after fusion of carbon in massive stars that will later become Type II supernovae: this fusion creates 26Mg, which, upon capturing free protons and neutrons becomes aluminium. Some smaller quantities of 27Al are created in hydrogen burning shells of evolved stars, where 26Mg can capture free protons.[48] Essentially all aluminium now in existence is 27Al; 26Al was present in the early Solar System but is currently extinct. However, the trace quantities of 26Al that do exist are the most common gamma ray emitter in the interstellar gas.[48]
+
+Overall, the Earth is about 1.59% aluminium by mass (seventh in abundance by mass).[49] Aluminium occurs in greater proportion in the Earth than in the Universe because aluminium easily forms the oxide and becomes bound into rocks and aluminium stays in the Earth's crust while less reactive metals sink to the core.[48] In the Earth's crust, aluminium is the most abundant (8.23% by mass[24]) metallic element and the third most abundant of all elements (after oxygen and silicon).[50] A large number of silicates in the Earth's crust contain aluminium.[51] In contrast, the Earth's mantle is only 2.38% aluminium by mass.[52] Aluminium also occurs in seawater at a concentration of 2 μg/kg.[24]
+
+Because of its strong affinity for oxygen, aluminium is almost never found in the elemental state; instead it is found in oxides or silicates. Feldspars, the most common group of minerals in the Earth's crust, are aluminosilicates. Aluminium also occurs in the minerals beryl, cryolite, garnet, spinel, and turquoise.[53] Impurities in Al2O3, such as chromium and iron, yield the gemstones ruby and sapphire, respectively.[54] Native aluminium metal can only be found as a minor phase in low oxygen fugacity environments, such as the interiors of certain volcanoes.[55] Native aluminium has been reported in cold seeps in the northeastern continental slope of the South China Sea. It is possible that these deposits resulted from bacterial reduction of tetrahydroxoaluminate Al(OH)4−.[56]
+
+Although aluminium is a common and widespread element, not all aluminium minerals are economically viable sources of the metal. Almost all metallic aluminium is produced from the ore bauxite (AlOx(OH)3–2x). Bauxite occurs as a weathering product of low iron and silica bedrock in tropical climatic conditions.[57] In 2017, most bauxite was mined in Australia, China, Guinea, and India.[58]
+
+The history of aluminium has been shaped by usage of alum. The first written record of alum, made by Greek historian Herodotus, dates back to the 5th century BCE.[59] The ancients are known to have used alum as a dyeing mordant and for city defense.[59] After the Crusades, alum, an indispensable good in the European fabric industry,[60] was a subject of international commerce;[61] it was imported to Europe from the eastern Mediterranean until the mid-15th century.[62]
+
+The nature of alum remained unknown. Around 1530, Swiss physician Paracelsus suggested alum was a salt of an earth of alum.[63] In 1595, German doctor and chemist Andreas Libavius experimentally confirmed this.[64] In 1722, German chemist Friedrich Hoffmann announced his belief that the base of alum was a distinct earth.[65] In 1754, German chemist Andreas Sigismund Marggraf synthesized alumina by boiling clay in sulfuric acid and subsequently adding potash.[65]
+
+Attempts to produce aluminium metal date back to 1760.[66] The first successful attempt, however, was completed in 1824 by Danish physicist and chemist Hans Christian Ørsted. He reacted anhydrous aluminium chloride with potassium amalgam, yielding a lump of metal looking similar to tin.[67][68][69] He presented his results and demonstrated a sample of the new metal in 1825.[70][71] In 1827, German chemist Friedrich Wöhler repeated Ørsted's experiments but did not identify any aluminium.[72] (The reason for this inconsistency was only discovered in 1921.)[73] He conducted a similar experiment in the same year by mixing anhydrous aluminium chloride with potassium and produced a powder of aluminium.[69] In 1845, he was able to produce small pieces of the metal and described some physical properties of this metal.[73] For many years thereafter, Wöhler was credited as the discoverer of aluminium.[74]
+
+As Wöhler's method could not yield great quantities of aluminium, the metal remained rare; its cost exceeded that of gold.[72] The first industrial production of aluminium was established in 1856 by French chemist Henri Etienne Sainte-Claire Deville and companions.[75] Deville had discovered that aluminium trichloride could be reduced by sodium, which was more convenient and less expensive than potassium, which Wöhler had used.[76] Even then, aluminium was still not of great purity and produced aluminium differed in properties by sample.[77]
+
+The first industrial large-scale production method was independently developed in 1886 by French engineer Paul Héroult and American engineer Charles Martin Hall; it is now known as the Hall–Héroult process.[78] The Hall–Héroult process converts alumina into the metal. Austrian chemist Carl Joseph Bayer discovered a way of purifying bauxite to yield alumina, now known as the Bayer process, in 1889.[79] Modern production of the aluminium metal is based on the Bayer and Hall–Héroult processes.[80]
+
+Prices of aluminium dropped and aluminium became widely used in jewelry, everyday items, eyeglass frames, optical instruments, tableware, and foil in the 1890s and early 20th century. Aluminium's ability to form hard yet light alloys with other metals provided the metal many uses at the time.[81] During World War I, major governments demanded large shipments of aluminium for light strong airframes.[82]
+
+By the mid-20th century, aluminium had become a part of everyday life and an essential component of housewares.[83] During the mid-20th century, aluminium emerged as a civil engineering material, with building applications in both basic construction and interior finish work,[84] and increasingly being used in military engineering, for both airplanes and land armor vehicle engines.[85] Earth's first artificial satellite, launched in 1957, consisted of two separate aluminium semi-spheres joined together and all subsequent space vehicles have used aluminium to some extent.[80] The aluminium can was invented in 1956 and employed as a storage for drinks in 1958.[86]
+
+Throughout the 20th century, the production of aluminium rose rapidly: while the world production of aluminium in 1900 was 6,800 metric tons, the annual production first exceeded 100,000 metric tons in 1916; 1,000,000 tons in 1941; 10,000,000 tons in 1971.[87] In the 1970s, the increased demand for aluminium made it an exchange commodity; it entered the London Metal Exchange, the oldest industrial metal exchange in the world, in 1978.[80] The output continued to grow: the annual production of aluminium exceeded 50,000,000 metric tons in 2013.[87]
+
+The real price for aluminium declined from $14,000 per metric ton in 1900 to $2,340 in 1948 (in 1998 United States dollars).[87] Extraction and processing costs were lowered over technological progress and the scale of the economies. However, the need to exploit lower-grade poorer quality deposits and the use of fast increasing input costs (above all, energy) increased the net cost of aluminium;[88] the real price began to grow in the 1970s with the rise of energy cost.[89] Production moved from the industrialized countries to countries where production was cheaper.[90] Production costs in the late 20th century changed because of advances in technology, lower energy prices, exchange rates of the United States dollar, and alumina prices.[91] The BRIC countries' combined share in primary production and primary consumption grew substantially in the first decade of the 21st century.[92] China is accumulating an especially large share of world's production thanks to abundance of resources, cheap energy, and governmental stimuli;[93] it also increased its consumption share from 2% in 1972 to 40% in 2010.[94] In the United States, Western Europe, and Japan, most aluminium was consumed in transportation, engineering, construction, and packaging.[95]
+
+Aluminium is named after alumina, or aluminium oxide in modern nomenclature. The word "alumina" comes from "alum", the mineral from which it was collected. The word "alum" comes from alumen, a Latin word meaning "bitter salt".[96] The word alumen stems from the Proto-Indo-European root *alu- meaning "bitter" or "beer".[97]
+
+British chemist Humphry Davy, who performed a number of experiments aimed to isolate the metal, is credited as the person who named the element. In 1808, he suggested the metal be named alumium in an article on his electrochemical research which was published in Philosophical Transactions of the Royal Society.[98] This suggestion was criticized by contemporary chemists from France, Germany, and Sweden, who insisted the metal should be named for the oxide, alumina, from which it would be isolated.[99] In 1812, Davy published a chemistry textbook in which he settled on the name aluminum, thus producing the modern name.[100] However, its spelling and pronunciation varies: aluminum is in use in the United States and Canada while aluminium is in use elsewhere.[101]
+
+The -ium suffix followed the precedent set in other newly discovered elements of the time: potassium, sodium, magnesium, calcium, and strontium (all of which Davy isolated himself). Nevertheless, element names ending in -um were known at the time, for example, platinum (known to Europeans since the 16th century), molybdenum (discovered in 1778), and tantalum (discovered in 1802). The -um suffix is consistent with the universal spelling alumina for the oxide (as opposed to aluminia); compare to lanthana, the oxide of lanthanum, and magnesia, ceria, and thoria, the oxides of magnesium, cerium, and thorium, respectively.
+
+In 1812, British scientist Thomas Young[102] wrote an anonymous review of Davy's book, in which he objected to aluminum and proposed the name aluminium: "for so we shall take the liberty of writing the word, in preference to aluminum, which has a less classical sound."[103] This name did catch on: while the -um spelling was occasionally used in Britain, the American scientific language used -ium from the start.[104] Most scientists used -ium throughout the world in the 19th century;[105] it still remains the standard in most other languages.[101] In 1828, American lexicographer Noah Webster used exclusively the aluminum spelling in his American Dictionary of the English Language.[106] In the 1830s, the -um spelling started to gain usage in the United States; by the 1860s, it had become the more common spelling there outside science.[104] In 1892, Hall used the -um spelling in his advertising handbill for his new electrolytic method of producing the metal, despite his constant use of the -ium spelling in all the patents he filed between 1886 and 1903. It was subsequently suggested this was a typo rather than intended.[101] By 1890, both spellings had been common in the U.S. overall, the -ium spelling being slightly more common; by 1895, the situation had reversed; by 1900, aluminum had become twice as common as aluminium; during the following decade, the -um spelling dominated American usage.[107] In 1925, the American Chemical Society adopted this spelling.[107]
+
+The International Union of Pure and Applied Chemistry (IUPAC) adopted aluminium as the standard international name for the element in 1990.[108] In 1993, they recognized aluminum as an acceptable variant;[108] the most recent 2005 edition of the IUPAC nomenclature of inorganic chemistry acknowledges this spelling as well.[109] IUPAC official publications use the -ium spelling as primary but list both where appropriate.[f]
+
+Aluminium production is highly energy-consuming, and so the producers tend to locate smelters in places where electric power is both plentiful and inexpensive.[112] As of 2012, the world's largest smelters of aluminium are located in China, Russia, Bahrain, United Arab Emirates, and South Africa.[113]
+
+In 2016, China was the top producer of aluminium with a world share of fifty-five percent; the next largest producing countries were Russia, Canada, India, and the United Arab Emirates.[111]
+
+According to the International Resource Panel's Metal Stocks in Society report, the global per capita stock of aluminium in use in society (i.e. in cars, buildings, electronics, etc.) is 80 kg (180 lb). Much of this is in more-developed countries (350–500 kg (770–1,100 lb) per capita) rather than less-developed countries (35 kg (77 lb) per capita).[114]
+
+Bauxite is converted to aluminium oxide by the Bayer process. Bauxite is blended for uniform composition and then is ground. The resulting slurry is mixed with a hot solution of sodium hydroxide; the mixture is then treated in a digester vessel at a pressure well above atmospheric, dissolving the aluminium hydroxide in bauxite while converting impurities into relatively insoluble compounds:[115]
+
+After this reaction, the slurry is at a temperature above its atmospheric boiling point. It is cooled by removing steam as pressure is reduced. The bauxite residue is separated from the solution and discarded. The solution, free of solids, is seeded with small crystals of aluminium hydroxide; this causes decomposition of the [Al(OH)4]− ions to aluminium hydroxide. After about half of aluminium has precipitated, the mixture is sent to classifiers. Small crystals of aluminium hydroxide are collected to serve as seeding agents; coarse particles are converted to aluminium oxide by heating; excess solution is removed by evaporation, (if needed) purified, and recycled.[115]
+
+The conversion of alumina to aluminium metal is achieved by the Hall–Héroult process. In this energy-intensive process, a solution of alumina in a molten (950 and 980 °C (1,740 and 1,800 °F)) mixture of cryolite (Na3AlF6) with calcium fluoride is electrolyzed to produce metallic aluminium. The liquid aluminium metal sinks to the bottom of the solution and is tapped off, and usually cast into large blocks called aluminium billets for further processing.[10]
+
+Anodes of the electrolysis cell are made of carbon—the most resistant material against fluoride corrosion—and either bake at the process or are prebaked. The former, also called Söderberg anodes, are less power-efficient and fumes released during baking are costly to collect, which is why they are being replaced by prebaked anodes even though they save the power, energy, and labor to prebake the cathodes. Carbon for anodes should be preferably pure so that neither aluminium nor the electrolyte is contaminated with ash. Despite carbon's resistivity against corrosion, it is still consumed at a rate of 0.4–0.5 kg per each kilogram of produced aluminium. Cathodes are made of anthracite; high purity for them is not required because impurities leach only very slowly. Cathode is consumed at a rate of 0.02–0.04 kg per each kilogram of produced aluminium. A cell is usually a terminated after 2–6 years following a failure of the cathode.[10]
+
+The Hall–Heroult process produces aluminium with a purity of above 99%. Further purification can be done by the Hoopes process. This process involves the electrolysis of molten aluminium with a sodium, barium, and aluminium fluoride electrolyte. The resulting aluminium has a purity of 99.99%.[10][116]
+
+Electric power represents about 20 to 40% of the cost of producing aluminium, depending on the location of the smelter. Aluminium production consumes roughly 5% of electricity generated in the United States.[108] Because of this, alternatives to the Hall–Héroult process have been researched, but none has turned out to be economically feasible.[10]
+
+Recovery of the metal through recycling has become an important task of the aluminium industry. Recycling was a low-profile activity until the late 1960s, when the growing use of aluminium beverage cans brought it to public awareness.[117] Recycling involves melting the scrap, a process that requires only 5% of the energy used to produce aluminium from ore, though a significant part (up to 15% of the input material) is lost as dross (ash-like oxide).[118] An aluminium stack melter produces significantly less dross, with values reported below 1%.[119]
+
+White dross from primary aluminium production and from secondary recycling operations still contains useful quantities of aluminium that can be extracted industrially. The process produces aluminium billets, together with a highly complex waste material. This waste is difficult to manage. It reacts with water, releasing a mixture of gases (including, among others, hydrogen, acetylene, and ammonia), which spontaneously ignites on contact with air;[120] contact with damp air results in the release of copious quantities of ammonia gas. Despite these difficulties, the waste is used as a filler in asphalt and concrete.[121]
+
+Aluminium is the most widely used non-ferrous metal.[122] The global production of aluminium in 2016 was 58.8 million metric tons. It exceeded that of any other metal except iron (1,231 million metric tons).[111]
+
+Aluminium is almost always alloyed, which markedly improves its mechanical properties, especially when tempered. For example, the common aluminium foils and beverage cans are alloys of 92% to 99% aluminium.[123] The main alloying agents are copper, zinc, magnesium, manganese, and silicon (e.g., duralumin) with the levels of other metals in a few percent by weight.[124]
+
+The major uses for aluminium metal are in:[125]
+
+The great majority (about 90%) of aluminium oxide is converted to metallic aluminium.[115] Being a very hard material (Mohs hardness 9),[126] alumina is widely used as an abrasive;[127] being extraordinarily chemically inert, it is useful in highly reactive environments such as high pressure sodium lamps.[128] Aluminium oxide is commonly used as a catalyst for industrial processes;[115] e.g. the Claus process to convert hydrogen sulfide to sulfur in refineries and to alkylate amines.[129][130] Many industrial catalysts are supported by alumina, meaning that the expensive catalyst material is dispersed over a surface of the inert alumina.[131] Another principal use is as a drying agent or absorbent.[115][132]
+
+Several sulfates of aluminium have industrial and commercial application. Aluminium sulfate (in its hydrate form) is produced on the annual scale of several millions of metric tons.[133] About two-thirds is consumed in water treatment.[133] The next major application is in the manufacture of paper.[133] It is also used as a mordant in dyeing, in pickling seeds, deodorizing of mineral oils, in leather tanning, and in production of other aluminium compounds.[133] Two kinds of alum, ammonium alum and potassium alum, were formerly used as mordants and in leather tanning, but their use has significantly declined following availability of high-purity aluminium sulfate.[133] Anhydrous aluminium chloride is used as a catalyst in chemical and petrochemical industries, the dyeing industry, and in synthesis of various inorganic and organic compounds.[133] Aluminium hydroxychlorides are used in purifying water, in the paper industry, and as antiperspirants.[133] Sodium aluminate is used in treating water and as an accelerator of solidification of cement.[133]
+
+Many aluminium compounds have niche applications, for example:
+
+Despite its widespread occurrence in the Earth's crust, aluminium has no known function in biology.[10] At pH 6–9 (relevant for most natural waters), aluminium precipitates out of water as the hydroxide and is hence not available; most elements behaving this way have no biological role or are toxic.[145] Aluminium salts are remarkably nontoxic, aluminium sulfate having an LD50 of 6207 mg/kg (oral, mouse), which corresponds to 500 grams for an 80 kg (180 lb) person.[10]
+
+In most people, aluminium is not as toxic as heavy metals. Aluminium is classified as a non-carcinogen by the United States Department of Health and Human Services.[146] There is little evidence that normal exposure to aluminium presents a risk to healthy adult,[147] and there is evidence of no toxicity if it is consumed in amounts not greater than 40 mg/day per kg of body mass.[146] Most aluminium consumed will leave the body in feces; most of the small part of it that enters the bloodstream, will be excreted via urine.[148]
+
+Aluminium, although rarely, can cause vitamin D-resistant osteomalacia, erythropoietin-resistant microcytic anemia, and central nervous system alterations. People with kidney insufficiency are especially at a risk.[146] Chronic ingestion of hydrated aluminium silicates (for excess gastric acidity control) may result in aluminium binding to intestinal contents and increased elimination of other metals, such as iron or zinc; sufficiently high doses (>50 g/day) can cause anemia.[146]
+
+During the 1988 Camelford water pollution incident people in Camelford had their drinking water contaminated with aluminium sulfate for several weeks. A final report into the incident in 2013 concluded it was unlikely that this had caused long-term health problems.[149]
+
+Aluminium has been suspected of being a possible cause of Alzheimer's disease,[150] but research into this for over 40 years has found, as of 2018[update], no good evidence of causal effect.[151][152]
+
+Aluminium increases estrogen-related gene expression in human breast cancer cells cultured in the laboratory.[153] In very high doses, aluminium is associated with altered function of the blood–brain barrier.[154] A small percentage of people[155] have contact allergies to aluminium and experience itchy red rashes, headache, muscle pain, joint pain, poor memory, insomnia, depression, asthma, irritable bowel syndrome, or other symptoms upon contact with products containing aluminium.[156]
+
+Exposure to powdered aluminium or aluminium welding fumes can cause pulmonary fibrosis.[157] Fine aluminium powder can ignite or explode, posing another workplace hazard.[158][159]
+
+Food is the main source of aluminium. Drinking water contains more aluminium than solid food;[146] however, aluminium in food may be absorbed more than aluminium from water.[160] Major sources of human oral exposure to aluminium include food (due to its use in food additives, food and beverage packaging, and cooking utensils), drinking water (due to its use in municipal water treatment), and aluminium-containing medications (particularly antacid/antiulcer and buffered aspirin formulations).[161] Dietary exposure in Europeans averages to 0.2–1.5 mg/kg/week but can be as high as 2.3 mg/kg/week.[146] Higher exposure levels of aluminium are mostly limited to miners, aluminium production workers, and dialysis patients.[162]
+
+Consumption of antacids, antiperspirants, vaccines, and cosmetics provide possible routes of exposure.[163] Consumption of acidic foods or liquids with aluminium enhances aluminium absorption,[164] and maltol has been shown to increase the accumulation of aluminium in nerve and bone tissues.[165]
+
+In case of suspected sudden intake of a large amount of aluminium, the only treatment is deferoxamine mesylate which may be given to help eliminate aluminium from the body by chelation.[166][167] However, this should be applied with caution as this reduces not only aluminium body levels, but also those of other metals such as copper or iron.[166]
+
+High levels of aluminium occur near mining sites; small amounts of aluminium are released to the environment at the coal-fired power plants or incinerators.[168] Aluminium in the air is washed out by the rain or normally settles down but small particles of aluminium remain in the air for a long time.[168]
+
+Acidic precipitation is the main natural factor to mobilize aluminium from natural sources[146] and the main reason for the environmental effects of aluminium;[169] however, the main factor of presence of aluminium in salt and freshwater are the industrial processes that also release aluminium into air.[146]
+
+In water, aluminium acts as a toxiс agent on gill-breathing animals such as fish by causing loss of plasma- and hemolymph ions leading to osmoregulatory failure.[169] Organic complexes of aluminium may be easily absorbed and interfere with metabolism in mammals and birds, even though this rarely happens in practice.[169]
+
+Aluminium is primary among the factors that reduce plant growth on acidic soils. Although it is generally harmless to plant growth in pH-neutral soils, in acid soils the concentration of toxic Al3+ cations increases and disturbs root growth and function.[170][171][172][173] Wheat has developed a tolerance to aluminium, releasing organic compounds that bind to harmful aluminium cations. Sorghum is believed to have the same tolerance mechanism.[174]
+
+Aluminium production possesses its own challenges to the environment on each step of the production process. The major challenge is the greenhouse gas emissions.[162] These gases result from electrical consumption of the smelters and the byproducts of processing. The most potent of these gases are perfluorocarbons from the smelting process.[162] Released sulfur dioxide is one of the primary precursors of acid rain.[162]
+
+A Spanish scientific report from 2001 claimed that the fungus Geotrichum candidum consumes the aluminium in compact discs.[175][176] Other reports all refer back to that report and there is no supporting original research. Better documented, the bacterium Pseudomonas aeruginosa and the fungus Cladosporium resinae are commonly detected in aircraft fuel tanks that use kerosene-based fuels (not avgas), and laboratory cultures can degrade aluminium.[177] However, these life forms do not directly attack or consume the aluminium; rather, the metal is corroded by microbe waste products.[178]

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+
+The Italy national football team (Italian: Nazionale di calcio dell'Italia) has officially represented Italy in international football since their first match in 1910. The squad is under the global jurisdiction of FIFA and is governed in Europe by UEFA—the latter of which was co-founded by the Italian team's supervising body, the Italian Football Federation (FIGC). Italy's home matches are played at various stadiums throughout Italy, and their primary training ground, Centro Tecnico Federale di Coverciano, is located at the FIGC technical headquarters in Coverciano, Florence.
+
+Italy is one of the most successful national teams in the history of the World Cup, having won four titles (1934, 1938, 1982, 2006) and appearing in two other finals (1970, 1994), reaching a third place (1990) and a fourth place (1978). In 1938, they became the first team to defend their World Cup title, and due to the outbreak of World War II, retained the title for a further 12 years. Italy had also previously won two Central European International Cups (1927–30, 1933–35). Between its first two World Cup victories, Italy won the Olympic football tournament (1936). After the majority of the team was killed in a plane crash in 1949, the team did not advance past the group stage of the following two World Cup tournaments, and also failed to qualify for the 1958 edition—failure to qualify for the World Cup would not happen again until the 2018 edition. Italy returned to form by 1968, winning a European Championship (1968), and after a period of alternating unsuccessful qualification rounds in Europe, later appeared in two other finals (2000, 2012). Italy's highest finish at the FIFA Confederations Cup was in 2013, where the squad achieved a third-place finish.
+
+The team is known as gli Azzurri (the Blues). Savoy blue is the common colour of the national teams representing Italy, as it is the traditional paint of the royal House of Savoy, which reigned over the Kingdom of Italy from 1860 to 1946. The national team is also known for its long-standing rivalries with other top footballing nations, such as those with Brazil, Croatia, France, Germany and Spain. In the FIFA World Rankings, in force since August 1993, Italy has occupied the first place several times, in November 1993 and during 2007 (February, April–June, September), with its worst placement in August 2018 in 21st place.
+
+The team's first match was held in Milan on 15 May 1910. Italy defeated France by a score of 6–2, with Italy's first goal scored by Pietro Lana.[3][4][5] The Italian team played with a (2–3–5) system and consisted of: De Simoni; Varisco, Calì; Trerè, Fossati, Capello; Debernardi, Rizzi, Cevenini I, Lana, Boiocchi. First captain of the team was Francesco Calì.[6]
+
+The first success in an official tournament came with the bronze medal in 1928 Summer Olympics, held in Amsterdam. After losing the semi-final against Uruguay, an 11–3 victory against Egypt secured third place in the competition. In the 1927–30 and 1933–35 Central European International Cup, Italy achieved the first place out of five Central European teams, topping the group with 11 points in both editions of the tournament.[7][8] Italy would also later win the gold medal at the 1936 Summer Olympics with a 2–1 victory in extra time in the gold medal match over Austria on 15 August 1936.[9]
+
+After declining to participate in the inaugural World Cup (1930, in Uruguay) the Italy national team won two consecutive editions of the tournament in 1934 and 1938, under the direction of coach Vittorio Pozzo and the performance of Giuseppe Meazza, who is considered one of the best Italian football players of all time by some.[10][11] Italy hosted the 1934 World Cup, and played their first ever World Cup match in a 7–1 win over the United States in Rome. Italy defeated Czechoslovakia 2–1 in extra time in the final in Rome, with goals by Raimundo Orsi and Angelo Schiavio to achieve their first World cup title in 1934. They achieved their second title in 1938 in a 4–2 defeat of Hungary, with two goals by Gino Colaussi and two goals by Silvio Piola in the World Cup that followed. Rumour has it, before the 1938 finals fascist Italian Prime Minister Benito Mussolini was to have sent a telegram to the team, saying "Vincere o morire!" (literally translated as "Win or die!"). However, no record remains of such a telegram, and World Cup player Pietro Rava said, when interviewed, "No, no, no, that's not true. He sent a telegram wishing us well, but no never 'win or die'."[12]
+
+In 1949, 10 of the 11 players in the team's initial line-up were killed in a plane crash that affected Torino, winners of the previous five Serie A titles. Italy did not advance further than the first round of the 1950 World Cup, as they were weakened severely due to the air disaster. The team had travelled by boat rather than by plane, fearing another accident.[13]
+
+In the World Cup finals of 1954 and 1962, Italy failed to progress past the first round, and did not qualify for the 1958 World Cup due to a 2–1 defeat to Northern Ireland in the last match of the qualifying round. Italy did not take part in the first edition of the European Championship in 1960 (then known as the European Nations Cup), and was knocked out by the Soviet Union in the first round of the 1964 European Nations' Cup qualifying.
+
+Their participation in the 1966 World Cup was ended by a 0–1 defeat at the hands of North Korea. Despite being the tournament favourites, the Azzurri, whose 1966 squad included Gianni Rivera and Giacomo Bulgarelli, were eliminated in the first round by the semi-professional North Koreans. The Italian team was bitterly condemned upon their return home, while North Korean scorer Pak Doo-ik was celebrated as the David who killed Goliath. Upon Italy's return home, furious fans threw fruit and rotten tomatoes at their transport bus at the airport.[14][15]
+
+In 1968, Italy participated in their first European Championship, hosting the European Championship and winning their first major competition since the 1938 World Cup, beating Yugoslavia in Rome for the title. The match holds the distinction of being the only European Championship or World Cup final to go to a replay.[16] After extra time the final ended in a 1–1 draw, and in the days before penalty shootouts, the rules required the match to be replayed a few days later. Italy won the replay 2–0 (with goals from Luigi Riva and Pietro Anastasi) to take the trophy.
+
+In the 1970 World Cup, exploiting the performances of European champions' players like Giacinto Facchetti, Gianni Rivera and Luigi Riva and with a new centre-forward Roberto Boninsegna, the team were able to come back to a World Cup final match after 32 years. They reached this result after one of the most famous matches in football history—the "Game of the Century", the 1970 World Cup semifinal between Italy and Germany that Italy won 4–3 in extra time, with five of the seven goals coming in extra time.[17] They were later defeated by Brazil in the final 4–1. The cycle of international successes ended in the 1974 World Cup, when the team was eliminated by Grzegorz Lato's Polish team in the first round.
+
+In the 1978 FIFA World Cup in Argentina, a new generation of Italian players, the most famous being Paolo Rossi, came to the international stage. Italy was the only team in the tournament to beat the eventual champions and host team Argentina. Second-round games against West Germany (0–0), Austria (1–0) and Netherlands (1–2) led Italy to the third-place final, where the team was defeated by Brazil 2–1. In the match that eliminated Italy from the tournament against the Netherlands, Italian goalkeeper Dino Zoff was beaten by a long-distance shot from Arie Haan, and Zoff was criticized for the defeat.[18] Italy hosted the 1980 UEFA European Football Championship, the first edition to be held between eight teams instead of four,[19] automatically qualifying for the finals as hosts. After two draws with Spain and Belgium and a narrow 1–0 win over England, Italy were beaten by Czechoslovakia in the third-place match on penalties 9–8 after Fulvio Collovati missed his kick.
+
+After a scandal in Serie A where some National team players such as Paolo Rossi[20] were prosecuted and suspended for match fixing and illegal betting, the Azzurri qualified for the second round of the 1982 World Cup after three uninspiring draws against Poland, Peru and Cameroon. Having been loudly criticized, the Italian team decided on a press black-out from then on, with only coach Enzo Bearzot and captain Dino Zoff appointed to speak to the press.
+
+Italy's regrouped in the second round group, a group of death with Argentina and Brazil. In the opener, Italy prevailed 2–1 over Argentina, with Italy's goals, both left-footed strikes, were scored by Marco Tardelli and Antonio Cabrini. After Brazil defeated Argentina 3–1, Italy needed to win in order to advance to the semi-finals. Twice Italy went in the lead with Paolo Rossi's goals, and twice Brazil came back. When Falcão scored to make it 2–2, Brazil would have been through on goal difference, but in the 74th minute Rossi scored the winning goal, for a hat-trick, in a crowded penalty area to send Italy to the semifinals after one of the greatest games in World Cup history.[21][22][23] Italy then progressed to the semi-final where they defeated Poland with two goals from Rossi.
+
+In the final, Italy met West Germany, who had advanced by a penalty shootout victory against France. The first half ended scoreless, after Antonio Cabrini missed a penalty awarded for a Hans-Peter Briegel foul on Bruno Conti. In the second half Paolo Rossi again scored the first goal, and while the Germans were pushing forward in search of an equaliser, Marco Tardelli and substitute Alessandro Altobelli finalised two contropiede counterattacks to make it 3–0. Paul Breitner scored home West Germany's consolation goal seven minutes from the end.[24]
+
+Tardelli's cry, "Gol! Gol!" was one of the defining images of Italy's 1982 World Cup triumph.[25] Paolo Rossi won the Golden Boot with six goals as well as the Golden Ball Award for the best player of the tournament,[26] and 40-year-old captain-goalkeeper Dino Zoff became the oldest player to win the World Cup.[27]
+
+However, Italy failed to qualify for the 1984 European Championship.[28][29] Italy then entered as reigning champions in the 1986 World Cup[30][31][32] but were eliminated by reigning European Champions, France, in the round of 16.[33]
+
+In 1986, Azeglio Vicini was appointed as new head coach, replacing Bearzot.[34] New coach conceded a chance to young players, such as Ciro Ferrara and Gianluca Vialli:[35] Sampdoria striker scored goals that gave Italy 1988 European Championship pass.[36] He was also shown like Altobelli's possibly successor, having his same goal attitude.[37] Both forwards stroke the target in Germany, where Soviet Union defeated azzurri in semi-finals.[38]
+
+Italy hosted the World Cup for the second time in 1990. The Italian attack featured talented forwards Salvatore Schillaci and a young Roberto Baggio. Italy played nearly all of their matches in Rome and did not concede a single goal in their first five matches, however, Italy lost in the semi-final to defending champion Argentina in Naples, losing 4–3 on penalty kicks following a 1–1 draw after extra time. Schillaci's first-half opener was equalised in the second half by Claudio Caniggia's header for Argentina. Aldo Serena missed the final penalty kick (with Roberto Donadoni also having his penalty saved by goalkeeper Sergio Goycochea). Italy went on to defeat England 2–1 in the third-place match in Bari, with Schillaci scoring the winning goal on a penalty to become the tournament's top scorer with six goals. Italy then failed to qualify for the 1992 European Championship. In November 1993, FIFA ranked Italy first in the FIFA World Rankings for their first time since the ranking system was introduced in December 1992.[39]
+
+At the 1994 World Cup in the United States, Italy lost the opening match against Ireland 0–1 at the Giants Stadium near New York City. After a 1–0 win against Norway in New York City and a 1–1 draw with Mexico at the RFK Stadium in Washington, D.C., Italy advanced from Group E based on goals scored among the four teams tied on points. During their round of 16 match at the Foxboro Stadium near Boston, Italy was down 0–1 late against Nigeria, but Baggio rescued Italy with an equaliser in the 88th minute and a penalty in extra time to take the win.[40] Baggio scored another late goal against Spain at their quarter-final match in Boston to seal a 2–1 win and two goals against Bulgaria in their semi-final match in New York City for another 2–1 win.[41][42]
+
+In the final, which took place in Los Angeles's Rose Bowl stadium 2,700 miles (4,320 km) and three time zones away from the Atlantic Northeast part of the United States where they had played all their previous matches, Italy, who had 24 hours less rest than Brazil, played 120 minutes of scoreless football, taking the match to a penalty shootout, the first time a World Cup final was settled in a penalty shootout.[43] Italy lost the subsequent shootout 3–2 after Baggio, who had been playing with the aid of a pain-killer injection[44] and a heavily bandaged hamstring,[45][46] missed the final penalty kick of the match, shooting over the crossbar.[47][48]
+
+Italy did not progress beyond the group stage at the finals of Euro 1996. Having defeated Russia 2–1 but losing to the Czech Republic by the same score, Italy required a win to be sure of progressing. Gianfranco Zola failed to convert a decisive penalty in a 0–0 draw against Germany,[49] who eventually won the tournament. During the qualifying campaign for the 1998 World Cup, Italy drew 0–0 to England on the last day of Group 2 matches as Italy finished in second place, one point behind England. Italy were then required to go through the play-off against Russia, advancing 2–1 on aggregate on 15 November 1997 with the winner coming from Pierluigi Casiraghi.[50] In the final tournament, Italy found themselves in another critical shootout for the third World Cup in a row. The Italian side, where Alessandro Del Piero and Baggio renewed the controversial staffetta ("relay") between Mazzola and Rivera from 1970, held the eventual World Champions and host team France to a 0–0 draw after extra time in the quarter-finals, but lost 4–3 in the shootout. With two goals scored in this tournament, Baggio is still the only Italian player to have scored in three different FIFA World Cup editions.[51]
+
+In the Euro 2000, another shootout decided Italy's fate but this time in their favour when defeating the co-hosts the Netherlands in the semi final. Italian goalkeeper Francesco Toldo saved one penalty during the match and two in the shootout, while the Dutch players missed one other penalty during the match and one during the shootout with a rate of one penalty scored out of six attempts. Emerging star Francesco Totti scored his penalty with a cucchiaio ("spoon") chip. Italy finished the tournament as runners-up, losing the final 2–1 against France (to a golden goal in extra time) after conceding les Bleus equalising goal just 30 seconds before the expected end of injury time (93rd minute).[52] After the defeat, coach Dino Zoff resigned in protest after being criticized by Milan club president and politician Silvio Berlusconi.[53]
+
+In the 2002 World Cup, a 2–0 victory against Ecuador with two Christian Vieri goals was followed by a series of controversial matches. During the match against Croatia, two goals were disallowed resulting in a 2–1 defeat for Italy. Despite two goals being ruled for borderline offsides, a late headed goal from Alessandro Del Piero helped Italy to a 1–1 draw with Mexico proving enough to advance to the knockout stages. However, co-host country South Korea eliminated Italy in the round of 16 by a score of 2–1. The game was highly controversial with members of the Italian team, most notably striker Francesco Totti and coach Giovanni Trapattoni, suggesting a conspiracy to eliminate Italy from the competition.[54] Trapattoni even obliquely accused FIFA of ordering the official to ensure a Korean victory so that one of the two host nations would remain in the tournament.[55] The most contentious decisions by the game referee Byron Moreno were an early penalty awarded to South Korea (saved by Buffon), a golden goal by Damiano Tommasi ruled offside, and the sending off of Totti after being presented with a second yellow card for an alleged dive in the penalty area.[56] FIFA President Sepp Blatter stated that the linesmen had been a "disaster" and admitted that Italy suffered from bad offside calls during the group matches, but he denied conspiracy allegations. While questioning Totti's sending off by Moreno, Blatter refused to blame Italy's loss entirely on the referees, stating: "Italy's elimination is not only down to referees and linesmen who made human not premeditated errors ... Italy made mistakes both in defense and in attack."[57]
+
+A three-way five point tie in the group stage of the 2004 European Championship left Italy as the "odd man out", as they failed to qualify for the quarter finals after finishing behind Denmark and Sweden on the basis of number of goals scored in matches among the tied teams. Italy's winning goal scored during stoppage time giving them a 2–1 victory over Bulgaria by Antonio Cassano proved futile, ending the team's tournament.
+
+The summer of 2004 marked the choice, by FIGC, to appoint Marcello Lippi for Italy's bench.[58] He made his debut in an upset 2–0 defeat in Iceland[59] but then managed to qualify for 2006 World Cup.[60][61] Italy's campaign in the tournament hosted by Germany was accompanied by open pessimism[62] due to the controversy caused by the 2006 Serie A scandal,[63] however these negative predictions were then refuted, as the Azzurri eventually won their fourth World Cup.
+
+Italy won their opening game against Ghana 2–0, with goals from Andrea Pirlo (40th minute) and substitute Vincenzo Iaquinta (83rd minute). The team performance was judged the best among the opening games by FIFA President Sepp Blatter.[64]
+
+The second match was a less convincing 1–1 draw with United States, with Alberto Gilardino's diving header equalized by a Cristian Zaccardo own goal.[65] After the equaliser, midfielder Daniele De Rossi and the United States's Pablo Mastroeni and Eddie Pope were sent off, leaving only nine men on the field for nearly the entirety of the second half, but the score remained unchanged despite a controversial decision when Gennaro Gattuso's shot was deflected in but disallowed because of an offside ruling. The same happened at the other end when U.S. winger DaMarcus Beasley's goal was not given due to teammate Brian McBride being ruled offside. De Rossi was suspended for four matches for elbowing McBride in the face and only returned for the final match.
+
+Italy finished first in Group E with a 2–0 win against the Czech Republic, with goals from defender Marco Materazzi (26th minute) and striker Filippo Inzaghi (87th minute), advancing to the Round of 16 in the knockout stages, where they faced Australia. In this match, Materazzi was controversially sent off early in the second half (53rd minute) after an attempted two-footed tackle on Australian midfielder Marco Bresciano. In stoppage time a controversial penalty kick was awarded to the Azzurri when referee Luis Medina Cantalejo ruled that Lucas Neill fouled Fabio Grosso. Francesco Totti converted into the upper corner of the goal past Mark Schwarzer for a 1–0 win.[66]
+
+In the quarter-finals, Italy beat Ukraine 3–0. Gianluca Zambrotta opened the scoring early (in the sixth minute) with a left-footed shot from outside the penalty area after a quick exchange with Totti created enough space. Luca Toni added two more goals in the second half (59th and 69th minute), as Ukraine pressed forward but were not able to score, hitting the crossbar and requiring several saves from Gianluigi Buffon and a goal-line clearance from Zambrotta. Afterwards, manager Marcello Lippi dedicated the victory to former Italian international Gianluca Pessotto, who was in the hospital recovering from an apparent suicide attempt.[67]
+
+In the semi-finals, Italy beat hosts Germany 2–0 with the two goals coming in the last two minutes of extra time. After a back-and-forth half-hour of extra time during which Alberto Gilardino and Gianluca Zambrotta struck the post and the crossbar respectively, Fabio Grosso scored in the 119th minute after a disguised Andrea Pirlo pass found him open in the penalty area for a bending left-footed shot into the far corner past German goalkeeper Jens Lehmann's dive. Substitute striker Alessandro Del Piero then sealed the victory by scoring with the last kick of the game at the end of a swift counterattack by Cannavaro, Totti and Gilardino.[68]
+
+The Azzurri won their fourth World Cup, defeating their long-time rivals France in Berlin, on 9 July, 5–3 on penalty kicks after a 1–1 draw at the end of extra time in the final. French captain Zinedine Zidane opened the scoring in the seventh minute with a chipped penalty kick, awarded for a controversial foul by Materazzi on Florent Malouda. Twelve minutes later, a header by Materazzi from a corner kick by Pirlo brought Italy even. In the second half, a potential winning goal by Toni was disallowed for a very close offside call by linesman Luc La Rossa. In the 110th minute, Zidane (playing in the last match of his career) was sent off by referee Horacio Elizondo for headbutting Materazzi in the chest after a verbal exchange;[69] Italy then won the penalty shootout 5–3, with the winner scored by Grosso; the crucial penalty miss from the French being David Trezeguet's, the same player who scored the golden goal for France in the Euro 2000. Trezeguet's attempt hit the crossbar, then shot down after its impact, and just stayed ahead of the line.[70]
+
+Ten different players scored for Italy in the tournament, and five goals out of twelve were scored by substitutes, while four goals were scored by defenders. Seven players — Gianluigi Buffon, Fabio Cannavaro, Gianluca Zambrotta, Andrea Pirlo, Gennaro Gattuso, Francesco Totti and Luca Toni — were named to the 23-man tournament All Star Team.[71] Buffon also won the Lev Yashin Award, given to the best goalkeeper of the tournament; he conceded only two goals in the tournament's seven matches, the first an own goal by Zaccardo and the second from Zidane's penalty kick in the final, and remained unbeaten for 460 consecutive minutes.[72]
+In honour of Italy winning the FIFA World Cup for a fourth time, all members of the World Cup-winning squad were awarded the Italian Order of Merit of Cavaliere Ufficiale.[73][74]
+
+Marcello Lippi, who had announced his resignation three days after the World Cup triumph, was replaced by Roberto Donadoni as the new coach of the Azzurri.[75] Italy played in the 2008 UEFA European Football Championship qualifying Group B, along with France. Italy won the group, with France being the runner-up. On 14 February 2007, Italy climbed to first in the FIFA World Rankings from second, with a total of 1,488 points, 37 points ahead of second ranked Argentina. This was the second time in the Azzurri's history that it had been ranked in first place, the first time being in 1993; they would also be ranked first several times throughout 2007, also in April–June and September.[39][76]
+
+In Euro 2008, the Azzurri lost 3–0 to the Netherlands. The following game against Romania ended 1–1, with a goal by Christian Panucci that came only one minute after Romania's Adrian Mutu capitalized on a mistake by Gianluca Zambrotta to give Romania the lead.[77] The result was preserved by Gianluigi Buffon who saved a penalty kick from Mutu in the 80th minute.[77]
+
+The final group game against France, a rematch of the 2006 World Cup Final, was a 2–0 Italy win. Andrea Pirlo scored from the penalty spot after a foul and red card for France defender Eric Abidal, and later a free kick by Daniele De Rossi took a deflection resulting Italy's second goal. Romania, entering the day a point ahead of the Italians in Group C, lost to the Netherlands 2–0, allowing Italy to pass into the quarter finals against eventual champions Spain, where they lost 2–4 on penalties after a 0–0 draw after 120 minutes. Within a week after the game, Roberto Donadoni's contract was terminated and Marcello Lippi was rehired as coach.[78]
+
+Italy qualified for their first ever FIFA Confederations Cup held in South Africa in June 2009 by virtue of winning the 2006 World Cup. They won their opening match of the tournament by a score of 3–1 against the United States, but subsequent defeats to Egypt (0–1) and Brazil (0–3) meant that they only finished third in the group on goals scored, and were eliminated.
+
+The national football team of Italy qualified for the 2010 FIFA World Cup after playing home games at Stadio Friuli, Stadio Via del Mare, Stadio San Nicola, Stadio Olimpico di Torino and Stadio Ennio Tardini. In October 2009, they achieved qualification after drawing with the Republic of Ireland 2–2. On 4 December 2009, the draw for the World Cup was made: Italy would be in Group F alongside three underdog teams: Paraguay, New Zealand and Slovakia.
+
+At the 2010 World Cup in South Africa, reigning champions Italy were unexpectedly eliminated in the first round, finishing last place in their group. After being held to 1–1 draws by Paraguay and New Zealand, they suffered a 3–2 loss to Slovakia.[79] It was the first time Italy failed to win a single game at a World Cup finals tournament, and in doing so became the third nation to be eliminated in the first round while holding the World Cup crown; the first being Brazil in 1966 and the second France in 2002.[80] Coincidentally, France who had been Italy's adversaries and the losing finalist in the 2006 World Cup, were also eliminated without winning a game in the first round in South Africa, making it the first time ever that neither finalist of the previous edition were able to reach the second round.[81]
+
+Marcello Lippi stepped down after Italy's World Cup campaign and was replaced by Cesare Prandelli, although Lippi's successor had already been announced before the tournament.[82] Italy began their campaign with Prandelli with a 1–0 loss to the Ivory Coast in a friendly match.[83] Then, during a Euro 2012 qualifier, Italy came back from behind to defeat Estonia 2–1. In the next Euro qualifier, Italy dominated the Faroe Islands 5–0. Italy then tied 0–0 with Northern Ireland. Five days later, Italy played Serbia; however, Serbian fans in Stadio Luigi Ferraris began to riot, throwing flares and shooting fireworks onto the pitch, subsequently causing the abandonment of the game.[84] Upon UEFA Disciplinary Review, Italy was awarded a 3–0 victory that propelled them to the top of their group.[85] In their first match of 2011, Italy drew 1–1 a friendly with Germany at Dortmund, in the same stadium where they beat Germany 2–0 to advance to the final of the 2006 World Cup. In March 2011, Italy won 1–0 over Slovenia to again secure its spot at the top of the qualification table. They then defeated Ukraine 2–0 in a friendly, despite being reduced to ten men for the late stages of the match. With their 3–0 defeat of Estonia in another Euro 2012 qualifier, Prandelli's Italy secured the table lead and also achieved 9 undefeated games in a row since their initial debacle. The streak was ended on 7 June 2011 by Trapattoni's current charges, the Republic of Ireland, with Italy losing 0–2 in a friendly in Liège.
+
+At the beginning of the second season under coach Prandelli, on 10 August 2011, Italy defeated the reigning world champions Spain for 2–1 in a friendly match played in Bari's Stadio San Nicola, but lost in a friendly to the United States, 1–0, on home soil on 29 February 2012.[86]
+
+Italy started their Euro 2012 campaign with a 1–1 draw to current reigning European and World champions Spain. In the following match, they draw 1–1 against Croatia. They finished second in their group behind Spain by beating the Republic of Ireland 2–0, which earned them a quarter-final match against the winners of group D, England. After a mostly one-sided affair in which Italy failed to take their chances, they managed to beat England on penalty kicks, even though they were down early in the shootout. A save by goalkeeper Gianluigi Buffon put them ahead after a chip shot from Andrea Pirlo. Prandelli's side won the shootout 4–2.[87][88]
+
+In their next game, the first semi-final of the competition, they faced Germany team who were tipped by many to be the next European champions.[89][90][91][92][93] However, two first-half goals by Mario Balotelli saw Germany sent home, and the Italians went through to the finals to face the title defenders Spain.
+
+In the final, however, they were unable to repeat their earlier performance against Spain, falling 4–0 to lose the championship. Prandelli's men were further undone by the string of injuries which left them playing with ten men for the last half-hour, as substitute Thiago Motta was forced to go off after all three substitutions had been made.[94]
+
+During the 2013 Confederations Cup in Brazil, Italy started in a group with Mexico, Japan and Brazil. After beating Mexico 2–1 and Japan 4–3, Italy eventually lost their final group game against tournament hosts Brazil 4–2. Italy then faced Spain in the semi-finals, in a rematch of the Euro 2012 final. Italy lost 7–6 (0–0 after extra time) in a penalty shoot-out after Leonardo Bonucci failed to score his kick.[95] Prandelli was praised for his tactics against the current World Cup and European champions.[96] Italy was then able to win the match for the third place by defeating Uruguay with the penalty score of 5–4 (2–2 after extra time).
+
+Italy was drawn in UEFA Group B for the 2014 World Cup qualification campaign. They won  the qualifying group without losing a match. Despite this successful run they were not seeded in pot 1 for the final seeding. In December 2013, Italy was drawn in Group D against Costa Rica, England and Uruguay. In its first match, Italy defeated England 2–1. However, in the second group stage match, underdogs Costa Rica beat the Italians 1–0.[97] In Italy's last group match, they were knocked out by Uruguay 1–0, due in part to two controversial calls from referee Marco Antonio Rodríguez (Mexico): in the 59th minute, midfielder Claudio Marchisio was sent off for a questionable tackle.[98] Later in the 80th minute, with the teams knotted at 0–0 which would have sent Italy to the next round, Uruguayan striker Luis Suárez bit defender Giorgio Chiellini on the shoulder but was not sent off.[99][100] Uruguay went on to score moments later in the 81st minute with a Diego Godín header from a corner kick, winning the game 1–0 and eliminating Italy. This marked Italy's second consecutive failure to reach the round of 16 at the World Cup finals. Shortly after this loss, coach Cesare Prandelli resigned.[101]
+
+The successful former Juventus manager Antonio Conte was selected to replace Cesare Prandelli as coach after the 2014 World Cup. Conte's debut as manager was against 2014 World Cup semi-finalists the Netherlands, in which Italy won 2–0. Italy's first defeat under Conte came ten games in to his empowerment from a 1–0 international friendly loss against Portugal on 16 June 2015.[102] On 10 October 2015, Italy qualified for Euro 2016, courtesy of a 3–1 win over Azerbaijan;[103] the result meant that Italy had managed to go 50 games unbeaten in European qualifiers.[104] Three days later, with a 2–1 win over Norway, Italy topped their Euro 2016 qualifying group with 24 points; four points clear of second placed Croatia.[105] However, with a similar fate to the 2014 World Cup group stage draw, Italy were not top seeded into the first pot. This had Italy see a draw with Belgium, Sweden and the Republic of Ireland in Group E.[106]
+
+On 4 April 2016, it was announced that Antonio Conte would step down as Italy coach after Euro 2016 to become head coach of English club Chelsea at the start of the 2016–17 Premier League season.[107] The 23-man squad, which was initially criticized by many fans and members of the media for its tactics and level of quality,[108] saw notable absences with Andrea Pirlo and Sebastian Giovinco controversially left out[109] and Claudio Marchisio and Marco Verratti omitted due to injury.[110][111] Italy opened Euro 2016 with a 2–0 victory over Belgium on 13 June.[112] Italy qualified for the round of 16 with one game to spare on 17 June with a lone goal by Éder for the victory against Sweden; the first time they won the second group game in a major international tournament since Euro 2000.[113] Italy also finished top of the group for the first time in a major tournament since the 2006 World Cup.[114] Italy defeated reigning European champions Spain 2–0 in the round of 16 match on 27 June.[115] Italy then faced off against the reigning World champions, rivals Germany, in the quarter-finals. Mesut Özil opened the scoring in the 65th minute for Germany, before Leonardo Bonucci converted a penalty in the 78th minute for Italy. The score remained 1–1 after extra time and Germany beat Italy 6–5 in the ensuing penalty shoot-out. It was the first time Germany overcame Italy in a major tournament, however, since the win occurred on penalties, it is statistically considered a draw.[116][117]
+
+For the 2018 FIFA World Cup qualification Italy were placed into the second pot due to being in 17th place in the FIFA World Rankings at the time of the group draws; Italy were drawn with Spain from pot one on 25 July 2015.[118] After Conte's planned departure following Euro 2016, Gian Piero Ventura took over as manager for the team, on 18 July 2016, signing a two-year contract.[119] His first match at the helm was a friendly against France, held at the Stadio San Nicola on 1 September, which ended in a 3–1 loss.[120] Four days later, he won his first competitive match in charge of Italy, the team's opening 2018 FIFA World Cup qualifier against Israel at Haifa, which ended in a 3–1 victory for Italy.[121]
+
+After Italy won all of their qualifying matches except for a 1–1 draw at home to Macedonia, as well as a 1–1 draw with Spain at home on 6 October 2016, and a 3–0 loss away to Spain on 2 September 2017, Italy finished in Group G in second place, five points behind Spain.[122][123] Italy were then required to go through the play-off against Sweden. After a 1–0 aggregate loss to Sweden, on 13 November 2017, Italy failed to qualify for the 2018 FIFA World Cup, the first time they failed to qualify for the World Cup since 1958.[124] Immediately following the match, veterans Giorgio Chiellini, Andrea Barzagli, Daniele De Rossi and captain Gianluigi Buffon all declared their retirement from the national team.[125][126][127][128][129] On 15 November 2017, Ventura was dismissed as head coach[130] and on 20 November 2017, Carlo Tavecchio resigned as president of the Italian Football Federation.[131][132]
+
+On 5 February 2018, the Italy U21 manager Luigi Di Biagio was appointed as the caretaker manager of the senior team.[133] On 17 March 2018, despite the initial decision to retire by veterans Buffon and Chiellini, they were both called up for Italy's March 2018 friendlies by caretaker manager Di Biagio.[134] Following the March friendlies against Argentina and England in which Italy were defeated and drew respectively, on 12 April 2018, Italy dropped six places to their lowest FIFA World Ranking at the time, to 20th place.[135] On 14 May 2018, Roberto Mancini was announced as the new manager.[136] On 28 May 2018, Italy won their first match under Mancini, a 2–1 victory in a friendly over Saudi Arabia.[137] On 16 August 2018, in the FIFA World Ranking that followed the 2018 World Cup, Italy dropped two places to their lowest ever ranking, to 21st place.[138] On 7 September 2018, Italy participated in the inaugural UEFA Nations League, drawing their first match of the tournament against Poland in Bologna with a score of 1–1.[139] On 12 October 2019, Italy qualified for Euro 2020 with three matches to spare after a 2–0 home win over Greece.[140] On 18 November, Italy finished Group J with ten wins in all ten of their matches, becoming only the sixth national side to qualify for a European Championship with a 100 per cent record, and the seventh instance, after France (1992 and 2004), Czech Republic (2000), Germany, Spain (both 2012) and England (2016).[141] On 17 March 2020, UEFA confirmed that Euro 2020 had been postponed by one year in response to the 2020 coronavirus pandemic in Europe.[142]
+
+The first shirt worn by the Italy national team, in its debut against France on 15 May 1910, was white. The choice of colour was due to the fact that a decision about the appearance of the kit had not yet been made, so it was decided not to have a colour, which was why white was chosen.[143] After two games, for a friendly against Hungary in Milan on 6 January 1911, the white shirt was replaced by a blue jersey (specifically savoy azure) — blue being the border colour of the royal House of Savoy crest used on the flag of the Kingdom of Italy (1861-1946); the shirt was accompanied by white shorts and black socks (which later became blue).[143] The team later became known as gli Azzurri (the Blues).[143][144][145]
+
+In the 1930s, Italy wore a black kit, ordered by the fascist regime of Benito Mussolini. The black kit debuted on 17 February 1935 in a friendly against France at the Stadio Nazionale PNF in Rome.[146] A blue shirt, white shorts and black socks were, however, worn at the 1936 Olympic Games in Berlin the following year. At the 1938 FIFA World Cup in France, the all-black kit was worn once (in the match against France).[147]
+
+After World War II, the fascist regime fell and the monarchy was abolished in 1946. The same year saw the birth of the Italian Republic, and the blue-and-white kit was reinstated. The cross of the former Royal House of Savoy was removed from the flag of Italy, and consequently from the national team's badge, now consisting solely of the Tricolore. For the 1954 FIFA World Cup, the country's name in Italian, "ITALIA", was placed above the tricolour shield, and for the 1982 FIFA World Cup, "FIGC", the abbreviation of the Italian Football Federation, was incorporated into the badge.[143]
+
+In 1983, to celebrate the victory at the World Cup of the previous year, three gold stars replaced the word "ITALIA" above the tricolour, representing their three World Cup victories until that point. And in 1984, a round emblem was launched, featuring the three stars, the inscriptions "ITALIA" and "FIGC", and the tricolour.[143]
+
+The first known kit manufacturer was Adidas in 1974. Since 2003, the kit has been made by Puma.[143] Since the 2000s, an all-blue uniform including blue shorts has occasionally been used, particularity in international tournaments.[143] After Italy's 2006 World Cup victory, a fourth star was added to the tricolour badge.
+
+Italy has five main rivalries with other top footballing nations.
+
+For the all-time record, see Italy national football team all-time record.
+
+Champions       Runners-up       Third place       Fourth place
+
+Minor titles:
+
+
+
+
+
+Current technical staff:[169]
+
+Giulio Nuciari  Fausto Salsano
+
+During the earliest days of Italian nation football, it was common for a Technical Commission to be appointed. The Commission took the role that a standard coach would currently play. Ever since 1967, the national team has been controlled only by the coach.
+
+For this reason, the coach of the Italy national team is still called Technical Commissioner (Commissario tecnico or CT, the use of this denomination has since then expanded into other team sports in Italy).
+
+Win
+  Draw
+  Loss
+  Fixtures
+
+The following players were called up for the UEFA Euro 2020 qualifying matches against Bosnia and Herzegovina on 15 November 2019 and Armenia on 18 November 2019.[170]
+
+Caps and goals updated as of 18 November 2019, after the match against Armenia.
+
+The following footballers have been selected in the past 12 months and are still eligible to represent.
+
+
+
+
+
+
+
+As of 18 November 2019, the players with the most appearances for Italy are:[171]
+
+Players in bold are still active in the national football team.
+
+As of 18 November 2019, the players with the most goals for Italy are:[172]
+
+Players in bold are still active in the national football team.
+
+List of captaincy periods of the various captains throughout the years.[173]
+
+For head to head records against other countries, see Italy national football team head to head.
+
+As of 18 November 2019, the complete official match record of the Italian national team comprises 823 matches: 437 wins, 224 draws and 162 losses.[183] During these matches, the team scored 1,433 times and conceded 818 goals. Italy's highest winning margin is nine goals, which has been achieved against the United States in 1948 (9–0). Their longest winning streak is 11 wins,[184] and their unbeaten record is 30 consecutive official matches.[185]

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+Desiderius Erasmus Roterodamus (/ˌdɛzɪˈdɪəriəs ɪˈræzməs/, English: Erasmus of Rotterdam;[note 1] 28 October 1466[3][4] – 12 July 1536) was a Dutch philosopher and Christian scholar who is widely considered to have been one of the greatest scholars of the northern Renaissance.[5] As a Catholic priest, Erasmus was an important figure in classical scholarship who wrote in a pure Latin style. Among humanists he enjoyed the sobriquet "Prince of the Humanists", and has been called "the crowning glory of the Christian humanists".[6] Using humanist techniques for working on texts, he prepared important new Latin and Greek editions of the New Testament, which raised questions that would be influential in the Protestant Reformation and Catholic Counter-Reformation. He also wrote On Free Will,[7] In Praise of Folly, Handbook of a Christian Knight, On Civility in Children, Copia: Foundations of the Abundant Style, Julius Exclusus, and many other works.
+
+Erasmus lived against the backdrop of the growing European religious Reformation. While he was critical of the abuses within the Catholic Church and called for reform, he nonetheless kept his distance from Luther, Henry VIII, and John Calvin and continued to recognise the authority of the pope, emphasizing a middle way with a deep respect for traditional faith, piety and grace, and rejecting Luther's emphasis on faith alone.[citation needed] Erasmus remained a member of the Catholic Church all his life, remaining committed to reforming the Church and its clerics' abuses from within.[8][9] He also held to the doctrine of synergism, which some Reformers (Calvinists) rejected in favor of the doctrine of monergism. His middle road ("via media") approach disappointed, and even angered, scholars in both camps.
+
+Erasmus died suddenly in Basel in 1536 while preparing to return to Brabant and was buried in Basel Minster, the former cathedral of the city.[10] A bronze statue of Erasmus was erected in 1622 in his city of birth, replacing an earlier work in stone.
+
+Desiderius Erasmus is reported to have been born in Rotterdam on 28 October in the late 1460s.[3][11] He was named after Saint Erasmus of Formiae, whom Erasmus's father Gerard personally favored.[12] A 17th-century legend has it that Erasmus was first named Geert Geerts (also Gerhard Gerhards or Gerrit Gerritsz),[13] but this is unfounded.[14] A well-known wooden picture indicates: Goudæ conceptus, Roterodami natus (Latin for Conceived in Gouda, born in Rotterdam). According to an article by historian Renier Snooy (1478–1537), Erasmus was born in Gouda.
+
+The exact year of his birth is controversial but most agree it was in 1466.[15] Evidence confirming the year of Erasmus' birth in 1466 can be found in his own words: fifteen out of twenty-three statements he made about his age indicate 1466.[16] He was christened "Erasmus" after the saint of that name.[17] Although associated closely with Rotterdam, he lived there for only four years, never to return. Information on his family and early life comes mainly from vague references in his writings. His parents were not legally married. His father, Gerard, was a Catholic priest and curate in Gouda.[18] Little is known of his mother, although her known name was Margaretha Rogerius (Latinized form of Dutch surname Rutgers)[19] and she was the daughter of a doctor from Zevenbergen. She may have been Gerard's housekeeper.[15][18][20] Although he was born out of wedlock, Erasmus was cared for by his parents until their early deaths from the plague in 1483. This solidified his view of his origin as a stain and cast a pall over his youth.[18]
+
+Erasmus was given the highest education available to a young man of his day, in a series of monastic or semi-monastic schools. At the age of nine, he and his older brother Peter were sent to one of the best Latin schools in the Netherlands, located at Deventer and owned by the chapter clergy of the Lebuïnuskerk (St Lebuin's Church),[15] though some earlier biographies assert it was a school run by the Brethren of the Common Life.[15] During his stay there the curriculum was renewed by the principal of the school, Alexander Hegius. For the first time ever in Europe, Greek was taught at a lower level than a university[dubious  – discuss] and this is where he began learning it.[21] He also gleaned there the importance of a personal relationship with God but eschewed the harsh rules and strict methods of the religious brothers and educators. His education there ended when plague struck the city about 1483, and his mother, who had moved to provide a home for her sons, died from the infection.[15]
+
+Most likely in 1487,[22] poverty[23] forced Erasmus into the consecrated life as a canon regular of St. Augustine at the canonry of Stein, in South Holland. He took vows there in late 1488[22] and was ordained to the Catholic priesthood on 25 April 1492.[23] It is said that he never seemed to have actively worked as a priest for a long time,[24] and certain abuses in religious orders were among the chief objects of his later calls to reform the Church from within.
+
+While at Stein, Erasmus met with a fellow canon, Servatius Rogerus,[25] and wrote a series of passionate letters in which he called Rogerus "half my soul".[26][27] This correspondence contrasts sharply with the generally detached and much more restrained attitude he showed in his later life. Later, while tutoring in Paris, he was suddenly dismissed by the guardian of Thomas Grey. Some have taken this as evidence of an illicit affair.[28] No personal denunciation was made of Erasmus during his lifetime, however, and he took pains in later life to distance these earlier episodes by condemning sodomy in his works, and praising sexual desire in marriage between men and women.[29]
+
+Soon after his priestly ordination, he got his chance to leave the canonry when offered the post of secretary to the Bishop of Cambrai, Henry of Bergen, on account of his great skill in Latin and his reputation as a man of letters. To allow him to accept that post, he was given a temporary dispensation from his religious vows on the grounds of poor health and love of Humanistic studies, though he remained a priest. Pope Leo X later made the dispensation permanent, a considerable privilege at the time.
+
+In 1495, with Bishop Henry's consent and a stipend, Erasmus went on to study at the University of Paris in the Collège de Montaigu, a centre of reforming zeal, under the direction of the ascetic Jan Standonck, of whose rigors he complained. The University was then the chief seat of Scholastic learning but already coming under the influence of Renaissance humanism. For instance, Erasmus became an intimate friend of an Italian humanist Publio Fausto Andrelini, poet and "professor of humanity" in Paris.
+
+The chief centres of Erasmus's activity were Paris, Leuven (in the Duchy of Brabant, now in Belgium), England, and Basel; yet he never belonged firmly in any one of these places. In 1499 he was invited to England by William Blount, 4th Baron Mountjoy, who offered to accompany him on his trip to England. According to Thomas Penn, Erasmus was "ever susceptible to the charms of attractive, well-connected, and rich young men".[30] His time in England was fruitful in the making of lifelong friendships with the leaders of English thought in the days of King Henry VIII: John Colet, Thomas More, John Fisher, Thomas Linacre and William Grocyn. At the University of Cambridge, he was the Lady Margaret's Professor of Divinity and had the option of spending the rest of his life as an English professor. He stayed at Queens' College, Cambridge, from 1510 to 1515.[31] His rooms were in the "I" staircase of Old Court, and he famously hated English ale and English weather.
+
+Erasmus suffered from poor health and complained that Queens' could not supply him with enough decent wine (wine was the Renaissance medicine for gallstones, from which Erasmus suffered). Until the early 20th century, Queens' College used to have a corkscrew that was purported to be "Erasmus' corkscrew", which was a third of a metre long; even today the college still has what it calls "Erasmus' chair".[32] Today Queens' College also has an Erasmus Building and an Erasmus Room. His legacy is marked for someone who complained bitterly about the lack of comforts and luxuries to which he was accustomed. As Queens' was an unusually humanist-leaning institution in the 16th century, Queens' College Old Library still houses many first editions of Erasmus' publications, many of which were acquired during that period by bequest or purchase, including Erasmus' New Testament translation, which is signed by friend and Polish religious reformer Jan Łaski.[33] Erasmus' friend, Chancellor John Fisher, was president of Queens' College from 1505 to 1508. His friendship with Fisher is the reason he chose to stay at Queens' while lecturing in Greek at the University.[34]
+
+During his first visit to England in 1499, he taught at the University of Oxford. Erasmus was particularly impressed by the Bible teaching of John Colet, who pursued a style more akin to the church fathers than the Scholastics. This prompted him, upon his return from England, to master the Greek language, which would enable him to study theology on a more profound level and to prepare a new edition of Jerome's Bible translation. On one occasion he wrote to Colet:
+
+I cannot tell you, dear Colet, how I hurry on, with all sails set, to holy literature. How I dislike everything that keeps me back, or retards me.[23]
+
+Despite a chronic shortage of money, he succeeded in learning Greek by an intensive, day-and-night study of three years, continuously begging in letters that his friends send him books and money for teachers.[35] Discovery in 1506 of Lorenzo Valla's New Testament Notes encouraged Erasmus to continue the study of the New Testament.[36]
+
+Erasmus preferred to live the life of an independent scholar and made a conscious effort to avoid any actions or formal ties that might inhibit his freedom of intellect and literary expression. Throughout his life, he was offered positions of honor and profit in academia but declined them all, preferring the uncertain but sufficient rewards of independent literary activity. He did however assist his friend John Colet by authoring Greek textbooks and procuring members of staff for the newly established St Paul's School.[37] From 1506 to 1509, he was in Italy: in 1506 he graduated as Doctor of Divinity at the Turin University, and he spent part of the time as a proofreader at the publishing house of Aldus Manutius in Venice. According to his letters, he was associated with the Venetian natural philosopher, Giulio Camillo,[38] but, apart from this, he had a less active association with Italian scholars than might have been expected.
+
+His residence at Leuven, where he lectured at the University, exposed Erasmus to much criticism from those ascetics, academics and clerics hostile to the principles of literary and religious reform and to the loose norms of the Renaissance adherents to which he was devoting his life. In 1517, he supported the foundation at the University, by his friend Hieronymus van Busleyden, of the Collegium Trilingue for the study of Hebrew, Latin, and Greek – after the model of the College of the Three Languages at the University of Alcalá. However, feeling that the lack of sympathy that prevailed at Leuven at that time was actually a form of mental persecution, he sought refuge in Basel, where under the shelter of Swiss hospitality he could express himself freely. Admirers from all quarters of Europe visited him there and he was surrounded by devoted friends, notably developing a lasting association with the great publisher Johann Froben.
+
+Only when he had mastered Latin did he begin to express himself on major contemporary themes in literature and religion. He felt called upon to use his learning in a purification of doctrine by returning to the historic documents and original languages of sacred Scripture. He tried to free the methods of scholarship from the rigidity and formalism of medieval traditions, but he was not satisfied with this. His revolt against certain forms of Christian monasticism and scholasticism was not based on doubts about the truth of doctrine, nor from hostility to the organization of the Church itself, nor from rejection of celibacy or monastic lifestyles. He saw himself as a preacher of righteousness by an appeal to reason, applied frankly and without fear of the magisterium. He always intended to remain faithful to Catholic doctrine and therefore was convinced he could criticize frankly virtually everyone and everything. Aloof from entangling obligations, Erasmus was the centre of the literary movement of his time, corresponding with more than five hundred men in the worlds of politics and of thought.
+
+In 1502, in Spain, Cardinal Francisco Jiménez de Cisneros had put together a team of Spanish translators to create a Compilation of a Bible in four languages, Greek, Hebrew, Aramaic and Latin. Translators for Greek were commissioned from Greece itself and worked closely with Latinists. Cardinal Cisneros's team completed and printed the full New Testament, including the Greek translation, in 1514. To do so they developed specific types to print Greek. Cisneros informed Erasmus of the works going on in Spain and may have sent a printed version of the New Testament to him. However, the Spanish team wanted the entire Bible to be released as one single work and withdrew from publication.
+
+The information and the delay allowed Erasmus to request a "Publication Privilege" of four years for the Greek New Testament to ensure that his work would be published first. He obtained it in 1516 from both Pope Leo X, to whom he would dedicate his work, and Emperor Maximilian I. Erasmus' Greek New Testament was published first, in 1516, forcing the Spanish team of Cisneros to wait until 1520 to publish their Complutensian Polyglot Bible.[39][40]
+
+It is hard to say if Erasmus' actions had an effect on delaying the publication of Complutensian Polyglot, causing the Spanish team to take more time, or if it made no difference in their perfectionism. The Spanish copy was approved for publication by the Pope in 1520; however, it was not released until 1522 due to the team's insistence on reviewing and editing. Only fifteen errors have been found in the entire six volumes and four languages of Cisneros's bible, an extraordinarily low number for the time. The fear of them publishing first, though, affected Erasmus's work, rushing him to printing and causing him to forego editing. The result was a large number of translation mistakes, transcription errors, and typos, that required further editions to be printed (see "publication").[41]
+
+Erasmus had been working for years on two projects: a collation of Greek texts and a fresh Latin New Testament. In 1512, he began his work on this Latin New Testament. He collected all the Vulgate manuscripts he could find to create a critical edition. Then he polished the language. He declared, "It is only fair that Paul should address the Romans in somewhat better Latin."[42] In the earlier phases of the project, he never mentioned a Greek text:
+
+My mind is so excited at the thought of emending Jerome’s text, with notes, that I seem to myself inspired by some god. I have already almost finished emending him by collating a large number of ancient manuscripts, and this I am doing at enormous personal expense.[43]
+
+While his intentions for publishing a fresh Latin translation are clear, it is less clear why he included the Greek text. Though some speculate that he intended to produce a critical Greek text or that he wanted to beat the Complutensian Polyglot into print, there is no evidence to support this. He wrote, "There remains the New Testament translated by me, with the Greek facing, and notes on it by me."[44] He further demonstrated the reason for the inclusion of the Greek text when defending his work:
+
+But one thing the facts cry out, and it can be clear, as they say, even to a blind man, that often through the translator’s clumsiness or inattention the Greek has been wrongly rendered; often the true and genuine reading has been corrupted by ignorant scribes, which we see happen every day, or altered by scribes who are half-taught and half-asleep.[45]
+
+So he included the Greek text to permit qualified readers to verify the quality of his Latin version. But by first calling the final product Novum Instrumentum omne ("All of the New Teaching") and later Novum Testamentum omne ("All of the New Testament") he also indicated clearly that he considered a text in which the Greek and the Latin versions were consistently comparable to be the essential core of the church's New Testament tradition.
+
+In a way it is legitimate to say that Erasmus "synchronized" or "unified" the Greek and the Latin traditions of the New Testament by producing an updated translation of both simultaneously. Both being part of canonical tradition, he clearly found it necessary to ensure that both were actually present in the same content. In modern terminology, he made the two traditions "compatible". This is clearly evidenced by the fact that his Greek text is not just the basis for his Latin translation, but also the other way round: there are numerous instances where he edits the Greek text to reflect his Latin version. For instance, since the last six verses of Revelation were missing from his Greek manuscript, Erasmus translated the Vulgate's text back into Greek. Erasmus also translated the Latin text into Greek wherever he found that the Greek text and the accompanying commentaries were mixed up, or where he simply preferred the Vulgate's reading to the Greek text.[46]
+
+Erasmus said it was "rushed into print rather than edited",[41] resulting in a number of transcription errors. After comparing what writings he could find, Erasmus wrote corrections between the lines of the manuscripts he was using (among which was Minuscule 2) and sent them as proofs to Froben.[47] His hurried effort was published by his friend Johann Froben of Basel in 1516 and thence became the first published Greek New Testament, the Novum Instrumentum omne, diligenter ab Erasmo Rot. Recognitum et Emendatum. Erasmus used several Greek manuscript sources because he did not have access to a single complete manuscript. Most of the manuscripts were, however, late Greek manuscripts of the Byzantine textual family and Erasmus used the oldest manuscript the least because "he was afraid of its supposedly erratic text."[48] He also ignored much older and better manuscripts that were at his disposal.[49]
+
+In the second (1519) edition, the more familiar term Testamentum was used instead of Instrumentum. This edition was used by Martin Luther in his German translation of the Bible, written for people who could not understand Latin. Together, the first and second editions sold 3,300 copies. By comparison, only 600 copies of the Complutensian Polyglot were ever printed. The first and second edition texts did not include the passage (1 John 5:7–8) that has become known as the Comma Johanneum. Erasmus had been unable to find those verses in any Greek manuscript, but one was supplied to him during production of the third edition.  The Catholic Church decreed that the Comma Johanneum was open to dispute (2 June 1927), and it is rarely included in modern scholarly translations.
+
+The third edition of 1522 was probably used by Tyndale for the first English New Testament (Worms, 1526) and was the basis for the 1550 Robert Stephanus edition used by the translators of the Geneva Bible and King James Version of the English Bible. Erasmus published a fourth edition in 1527 containing parallel columns of Greek, Latin Vulgate and Erasmus' Latin texts. In this edition Erasmus also supplied the Greek text of the last six verses of Revelation (which he had translated from Latin back into Greek in his first edition) from Cardinal Ximenez's Biblia Complutensis. In 1535 Erasmus published the fifth (and final) edition which dropped the Latin Vulgate column but was otherwise similar to the fourth edition. Later versions of the Greek New Testament by others, but based on Erasmus' Greek New Testament, became known as the Textus Receptus.[50]
+
+Erasmus dedicated his work to Pope Leo X as a patron of learning and regarded this work as his chief service to the cause of Christianity. Immediately afterwards, he began the publication of his Paraphrases of the New Testament, a popular presentation of the contents of the several books. These, like all of his writings, were published in Latin but were quickly translated into other languages with his encouragement.
+
+Erasmus, in his capacity as humanist editor, advised major printers such as Aldus Manutius on which manuscripts to publish.[51]
+
+The Protestant Reformation began in the year following the publication of his edition of the Greek New Testament (1516) and tested Erasmus' character. The issues between the Catholic Church and the growing religious movement which would later become known as Protestantism, had become so clear that few could escape the summons to join the debate. Erasmus, at the height of his literary fame, was inevitably called upon to take sides, but partisanship was foreign to his nature and his habits. Despite all his criticism of clerical corruption and abuses within the Catholic Church,[8] which lasted for years and was also directed towards many of the Church's basic doctrines,[9] Erasmus shunned the Reformation movement along with its most radical and reactionary offshoots,[8][need quotation to verify] and sided with neither party.[8]
+
+The world had laughed at his satire, but few had interfered with his activities. He believed that his work so far had commended itself to the best minds and also to the dominant powers in the religious world. Erasmus did not build a large body of supporters with his letters. He chose to write in Greek and Latin, the languages of scholars. His critiques reached an elite but small audience.[52]
+
+"Free will does not exist", according to Luther in his letter De Servo Arbitrio to Erasmus translated into German by Justus Jonas (1526), in that sin makes human beings completely incapable of bringing themselves to God. Noting Luther's criticism of the Catholic Church, Erasmus described him as "a mighty trumpet of gospel truth" while agreeing, "It is clear that many of the reforms for which Luther calls are urgently needed."[53] He had great respect for Luther, and Luther spoke with admiration of Erasmus' superior learning. Luther hoped for his cooperation in a work which seemed only the natural outcome of his own. In their early correspondence, Luther expressed boundless admiration for all Erasmus had done in the cause of a sound and reasonable Christianity and urged him to join the Lutheran party. Erasmus declined to commit himself, arguing that to do so would endanger his position as a leader in the movement for pure scholarship which he regarded as his purpose in life. Only as an independent scholar could he hope to influence the reform of religion. When Erasmus hesitated to support him, the straightforward Luther became angered that Erasmus was avoiding the responsibility due either to cowardice or a lack of purpose. However, any hesitancy on the part of Erasmus stemmed, not from lack of courage or conviction, but rather from a concern over the mounting disorder and violence of the reform movement. To Philip Melanchthon in 1524 he wrote:
+
+I know nothing of your church; at the very least it contains people who will, I fear, overturn the whole system and drive the princes into using force to restrain good men and bad alike. The gospel, the word of God, faith, Christ, and Holy Spirit – these words are always on their lips; look at their lives and they speak quite another language.[54]
+
+Again, in 1529, he writes “An epistle against those who falsely boast they are Evangelicals”[55] to Vulturius Neocomus (Gerardus Geldenhouwer). Here Erasmus complains of the doctrines and morals of the Reformers:
+
+You declaim bitterly against the luxury of priests, the ambition of bishops, the tyranny of the Roman Pontiff, and the babbling of the sophists; against our prayers, fasts, and Masses; and you are not content to retrench the abuses that may be in these things, but must needs abolish them entirely. ...Look around on this ‘Evangelical’ generation,[56] and observe whether amongst them less indulgence is given to luxury, lust, or avarice, than amongst those whom you so detest. Show me any one person who by that Gospel has been reclaimed from drunkenness to sobriety, from fury and passion to meekness, from avarice to liberality, from reviling to well-speaking, from wantonness to modesty. I will show you a great many who have become worse through following it. ...The solemn prayers of the Church are abolished, but now there are very many who never pray at all. ...I have never entered their conventicles, but I have sometimes seen them returning from their sermons, the countenances of all of them displaying rage, and wonderful ferocity, as though they were animated by the evil spirit. ...Who ever beheld in their meetings any one of them shedding tears, smiting his breast, or grieving for his sins? ...Confession to the priest is abolished, but very few now confess to God. ...They have fled from Judaism that they may become Epicureans.[57]
+
+Apart from these perceived moral failings of the Reformers, Erasmus also dreaded any change in doctrine, citing the long history of the Church as a bulwark against innovation. In book I of his Hyperaspistes he puts the matter bluntly to Luther:
+
+We are dealing with this: Would a stable mind depart from the opinion handed down by so many men famous for holiness and miracles, depart from the decisions of the Church, and commit our souls to the faith of someone like you who has sprung up just now with a few followers, although the leading men of your flock do not agree either with you or among themselves – indeed though you do not even agree with yourself, since in this same Assertion[58] you say one thing in the beginning and something else later on, recanting what you said before.[59]
+
+Continuing his chastisement of Luther – and undoubtedly put off by the notion of there being "no pure interpretation of Scripture anywhere but in Wittenberg"[60] – Erasmus touches upon another important point of the controversy:
+
+You stipulate that we should not ask for or accept anything but Holy Scripture, but you do it in such a way as to require that we permit you to be its sole interpreter, renouncing all others. Thus the victory will be yours if we allow you to be not the steward but the lord of Holy Scripture.[61]
+
+Though he sought to remain firmly neutral in doctrinal disputes, each side accused him of siding with the other, perhaps because of his neutrality. It was not for lack of fidelity with either side but a desire for fidelity with them both:
+
+I detest dissension because it goes both against the teachings of Christ and against a secret inclination of nature. I doubt that either side in the dispute can be suppressed without grave loss.[53]
+
+In his catechism (entitled Explanation of the Apostles' Creed) (1533), Erasmus took a stand against Luther's teaching by asserting the unwritten Sacred Tradition as just as valid a source of revelation as the Bible, by enumerating the Deuterocanonical books in the canon of the Bible and by acknowledging seven sacraments.[62] He called "blasphemers" anyone who questioned the perpetual virginity of Mary.[63] However, he supported lay access to the Bible.[63]
+
+In a letter to Nikolaus von Amsdorf, Luther objected to Erasmus’ catechism and called Erasmus a "viper," "liar," and "the very mouth and organ of Satan".[64]
+
+As regards the Reformation, Erasmus was accused by the monks to have:
+
+prepared the way and was responsible for Martin Luther. Erasmus, they said, had laid the egg, and Luther had hatched it. Erasmus wittily dismissed the charge, claiming that Luther had hatched a different bird entirely.[65]
+
+Twice in the course of the great discussion, he allowed himself to enter the field of doctrinal controversy, a field foreign to both his nature and his previous practice. One of the topics he dealt with was free will, a crucial question. In his De libero arbitrio diatribe sive collatio (1524), he lampoons the Lutheran view on free will. He lays down both sides of the argument impartially. The "Diatribe" did not encourage any definite action; this was its merit to the Erasmians and its fault in the eyes of the Lutherans. In response, Luther wrote his De servo arbitrio (On the Bondage of the Will, 1525), which attacks the "Diatribe" and Erasmus himself, going so far as to claim that Erasmus was not a Christian. Erasmus responded with a lengthy, two-part Hyperaspistes (1526–27). In this controversy Erasmus lets it be seen that he would like to claim more for free will than St. Paul and St. Augustine seem to allow according to Luther's interpretation.[66] For Erasmus the essential point is that humans have the freedom of choice.[67] The conclusions Erasmus reached drew upon a large array of notable authorities, including, from the Patristic period, Origen, John Chrysostom, Ambrose, Jerome, and Augustine, in addition to many leading Scholastic authors, such as Thomas Aquinas and Duns Scotus. The content of Erasmus' works also engaged with later thought on the state of the question, including the perspectives of the via moderna school and of Lorenzo Valla, whose ideas he rejected.
+
+As the popular response to Luther gathered momentum, the social disorders, which Erasmus dreaded and Luther disassociated himself from, began to appear, including the German Peasants' War, the Anabaptist disturbances in Germany and in the Low Countries, iconoclasm, and the radicalization of peasants across Europe. If these were the outcomes of reform, he was thankful that he had kept out of it. Yet he was ever more bitterly accused of having started the whole "tragedy" (as the Catholics dubbed Protestantism).
+
+When the city of Basel definitely adopted the Reformation in 1529, Erasmus gave up his residence there and settled in the imperial town of Freiburg im Breisgau.
+
+Certain works of Erasmus laid a foundation for religious toleration and ecumenism. For example, in De libero arbitrio, opposing certain views of Martin Luther, Erasmus noted that religious disputants should be temperate in their language, "because in this way the truth, which is often lost amidst too much wrangling may be more surely perceived." Gary Remer writes, "Like Cicero, Erasmus concludes that truth is furthered by a more harmonious relationship between interlocutors."[68] Although Erasmus did not oppose the punishment of heretics, in individual cases he generally argued for moderation and against the death penalty. He wrote, "It is better to cure a sick man than to kill him."[69]
+
+A test of the Reformation was the doctrine of the sacraments, and the crux of this question was the observance of the Eucharist. In 1530, Erasmus published a new edition of the orthodox treatise of Algerus against the heretic Berengar of Tours in the eleventh century. He added a dedication, affirming his belief in the reality of the Body of Christ after consecration in the Eucharist, commonly referred to as transubstantiation. The sacramentarians, headed by Œcolampadius of Basel, were, as Erasmus says, quoting him as holding views similar to their own in order to try to claim him for their schismatic and "erroneous" movement.
+
+When his strength began to fail, he decided to accept an invitation by Queen Mary of Hungary, Regent of the Netherlands, to move from Freiburg to Brabant. However, during preparations for the move in 1536, he suddenly died from an attack of dysentery during a visit to Basel.[70] He had remained loyal to the papal authorities in Rome, but he did not have the opportunity to receive the last rites of the Catholic Church; the reports of his death do not mention whether he asked for a priest or not. According to Jan van Herwaarden, this is consistent with his view that outward signs were not important; what mattered is the believer's direct relationship with God, which he noted "as the [Catholic] church believes". However, Herwaarden observes that "he did not dismiss the rites and sacraments out of hand but asserted a dying person could achieve a state of salvation without the priestly rites, provided their faith and spirit were attuned to God."[71]
+He was buried with great ceremony in Basel Minster (the former cathedral) there.[10]
+
+His last words, as recorded by his friend Beatus Rhenanus, were apparently "Dear God" (Dutch: Lieve God).[72] A bronze statue of him was erected in the city of his birth in 1622, replacing an earlier work in stone.
+
+Erasmus wrote both on church subjects and those of general human interest. By the 1530s, the writings of Erasmus accounted for 10 to 20 percent of all book sales in Europe.[73]  He is credited with coining the adage, "In the land of the blind, the one-eyed man is king."  With the collaboration of Publio Fausto Andrelini, he formed a paremiography (collection) of Latin proverbs and adages, commonly titled Adagia.  Erasmus is also generally credited with originating the phrase "Pandora's box", arising through an error in his translation of Hesiod's Pandora in which he confused pithos (storage jar) with pyxis (box).
+
+His more serious writings begin early with the Enchiridion militis Christiani, the "Handbook of the Christian Soldier" (1503 – translated into English a few years later by the young William Tyndale).  (A more literal translation of enchiridion – "dagger" – has been likened to "the spiritual equivalent of the modern Swiss Army knife.")[74] In this short work, Erasmus outlines the views of the normal Christian life, which he was to spend the rest of his days elaborating.  The chief evil of the day, he says, is formalism – going through the motions of tradition without understanding their basis in the teachings of Christ.  Forms can teach the soul how to worship God, or they may hide or quench the spirit. In his examination of the dangers of formalism, Erasmus discusses monasticism, saint worship, war, the spirit of class and the foibles of "society."
+
+The Enchiridion is more like a sermon than a satire. With it Erasmus challenged common assumptions, painting the clergy as educators who should share the treasury of their knowledge with the laity. He emphasized personal spiritual disciplines and called for a reformation which he characterized as a collective return to the Fathers and Scripture. Most importantly, he extolled the reading of scripture as vital because of its power to transform and motivate toward love. Much like the Brethren of the Common Life, he wrote that the New Testament is the law of Christ people are called to obey and that Christ is the example they are called to imitate.
+
+According to Ernest Barker, "Besides his work on the New Testament, Erasmus laboured also, and even more arduously, on the early Fathers. …Among the Latin Fathers he edited the works of St Jerome, St Hilary, and St Augustine; among the Greeks he worked on Irenaeus, Origen and Chrysostom."[75]
+
+Erasmus also wrote of the legendary Frisian freedom fighter and rebel Pier Gerlofs Donia (Greate Pier), though more often in criticism than in praise of his exploits.  Erasmus saw him as a dim, brutal man who preferred physical strength to wisdom.[76]
+
+One of Erasmus' best-known works, inspired by De triumpho stultitiae (written by Italian humanist Faustino Perisauli), is The Praise of Folly, published under the double title Moriae encomium (Greek, Latinised) and Laus stultitiae (Latin).[77] A satirical attack on superstitions and other traditions of European society in general and the western Church in particular, it was written in 1509, published in 1511, and dedicated to Sir Thomas More, whose name the title puns.
+
+The Institutio principis Christiani  or "Education of a Christian Prince" (Basel, 1516) was written as advice to the young king Charles of Spain (later Charles V, Holy Roman Emperor). Erasmus applies the general principles of honor and sincerity to the special functions of the Prince, whom he represents throughout as the servant of the people. Education was published in 1516, three years after[78] Niccolò Machiavelli’s The Prince was written; a comparison between the two is worth noting. Machiavelli stated that, to maintain control by political force, it is safer for a prince to be feared than loved. Erasmus preferred for the prince to be loved, and strongly suggested a well-rounded education in order to govern justly and benevolently and avoid becoming a source of oppression.
+
+As a result of his reformatory activities, Erasmus found himself at odds with both of the great parties. His last years were embittered by controversies with men toward whom he was sympathetic. Notable among these was Ulrich von Hutten, a brilliant but erratic genius who had thrown himself into the Lutheran cause and declared that Erasmus, if he had a spark of honesty, would do the same. In his reply in 1523, Spongia adversus aspergines Hutteni, Erasmus displays his skill in semantics. He accuses Hutten of having misinterpreted his utterances about reform and reiterates his determination never to break with the Church.
+
+The Ciceronianus came out in 1528, attacking the style of Latin that was based exclusively and fanatically on Cicero's writings. Etienne Dolet wrote a riposte titled Erasmianus in 1535.
+
+Erasmus' last major work, published the year of his death, is the Ecclesiastes or "Gospel Preacher" (Basel, 1536), a massive manual for preachers of around a thousand pages. Though somewhat unwieldy because Erasmus was unable to edit it properly in his old age, it is in some ways the culmination of all of Erasmus' literary and theological learning, offering prospective preachers advice on nearly every conceivable aspect of their vocation with extraordinarily abundant reference to classical and biblical sources.
+
+Erasmus’ Sileni Alcibiadis is one of his most direct assessments of the need for Church reform. Johann Froben published it first within a revised edition of the Adagia in 1515, then as a stand-alone work in 1517. This essay has been likened to John Colet's Convocation Sermon, though the styles differ.
+
+Sileni is the plural (Latin) form of Silenus, a creature often related to the Roman wine god Bacchus and represented in pictorial art as inebriated, merry revellers, variously mounted on donkeys, singing, dancing, playing flutes, etc. Alcibiades was a Greek politician in the 5th century BCE and a general in the Peloponnesian War; he figures here more as a character written into some of Plato's dialogues – a young, debauched playboy whom Socrates tries to convince to seek truth instead of pleasure, wisdom instead of pomp and splendor.
+
+The term Sileni – especially when juxtaposed with the character of Alcibiades – can therefore be understood as an evocation of the notion that something on the inside is more expressive of a person's character than what one sees on the outside. For instance, something or someone ugly on the outside can be beautiful on the inside, which is one of the main points of Plato's dialogues featuring Alcibiades and the Symposion, in which Alcibiades also appears.
+
+In support of this, Erasmus states, "Anyone who looks closely at the inward nature and essence will find that nobody is further from true wisdom than those people with their grand titles, learned bonnets, splendid sashes and bejeweled rings, who profess to be wisdom’s peak."  On the other hand, Erasmus lists several Sileni and then questions whether Christ is the most noticeable Silenus of them all. The Apostles were Sileni since they were ridiculed by others. He believes that the things which are the least ostentatious can be the most significant, and that the Church constitutes all Christian people – that despite contemporary references to clergy as the whole of the Church, they are merely its servants. He criticizes those that spend the Church's riches at the people's expense. The true point of the Church is to help people lead Christian lives. Priests are supposed to be pure, yet when they stray, no one condemns them. He criticizes the riches of the popes, believing that it would be better for the Gospel to be most important.
+
+The popularity of his books is reflected in the number of editions and translations that have appeared since the sixteenth century. Ten columns of the catalogue of the British Library are taken up with the enumeration of the works and their subsequent reprints. The greatest names of the classical and patristic world are among those translated, edited, or annotated by Erasmus, including Saint Ambrose, Aristotle, Saint Augustine, Saint Basil, Saint John Chrysostom, Cicero and Saint Jerome.[79]
+
+In his native Rotterdam, the University and Gymnasium Erasmianum have been named in his honor. In 2003, a poll showing that most Rotterdammers believed Erasmus to be the designer of the local "Erasmus Bridge" instigated the founding of the Erasmus House (Rotterdam),[80] and the Erasmus House (Jakarta)[81] dedicated to celebrating Erasmus' legacy. Three moments in Erasmus' life are celebrated annually. On 1 April, the city celebrates the publication of his best-known book The Praise of Folly. On 11 July, the Night of Erasmus celebrates the lasting influence of his work. His birthday is celebrated on 28 October.[82]
+
+Erasmus' reputation and the interpretations of his work have varied over time. Moderate Catholics recognized him as a leading figure in attempts to reform the Church, while Protestants recognized his initial support for Luther's ideas and the groundwork he laid for the future Reformation, especially in biblical scholarship. By the 1560s, however, there was a marked change in reception.[citation needed]
+
+According to Franz Anton Knittel, Erasmus in his Novum Instrumentum omne did not incorporate the Comma from the Codex Montfortianus (concerning the Trinity), because of grammar differences, but used Complutensian Polyglotta. According to him the Comma was known to Tertullian.[83]
+
+Protestant views of Erasmus fluctuated depending on region and period, with continual support in his native Netherlands and in cities of the Upper Rhine area. However, following his death and in the late sixteenth century, many Reformation supporters saw Erasmus's critiques of Luther and lifelong support for the universal Catholic Church as damning, and second-generation Protestants were less vocal in their debts to the great humanist. Nevertheless, his reception is demonstrable among Swiss Protestants in the sixteenth century: he had an indelible influence on the biblical commentaries of, for example, Konrad Pellikan, Heinrich Bullinger, and John Calvin, all of whom used both his annotations on the New Testament and his paraphrases of same in their own New Testament commentaries.[84]
+
+However, Erasmus designated his own legacy, and his life works were turned over at his death to his friend the Protestant humanist turned remonstrator Sebastian Castellio for the repair of the breach and divide of Christianity in its Catholic, Anabaptist, and Protestant branches.[85]
+
+By the coming of the Age of Enlightenment, however, Erasmus increasingly again became a more widely respected cultural symbol and was hailed as an important figure by increasingly broad groups. In a letter to a friend, Erasmus once had written: "That you are patriotic will be praised by many and easily forgiven by everyone; but in my opinion it is wiser to treat men and things as though we held this world the common fatherland of all."[86] Thus, the universalist ideals of Erasmus are sometimes claimed to be important for fixing global governance.[87]
+
+Several schools, faculties and universities in the Netherlands and Belgium are named after him, as is Erasmus Hall in Brooklyn, New York, USA.The European Union's Erasmus Programme scholarships enable students to spend up to a year of their university courses in a university in another European country.
+
+Eramus is credited with saying "When I get a little money I buy books; and if any is left, I buy food and clothes."[88]
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+He is also blamed for the mistranslation from Greek of to call a bowl a bowl as to call a spade a spade.[89]
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+Desiderius Erasmus Roterodamus (/ˌdɛzɪˈdɪəriəs ɪˈræzməs/, English: Erasmus of Rotterdam;[note 1] 28 October 1466[3][4] – 12 July 1536) was a Dutch philosopher and Christian scholar who is widely considered to have been one of the greatest scholars of the northern Renaissance.[5] As a Catholic priest, Erasmus was an important figure in classical scholarship who wrote in a pure Latin style. Among humanists he enjoyed the sobriquet "Prince of the Humanists", and has been called "the crowning glory of the Christian humanists".[6] Using humanist techniques for working on texts, he prepared important new Latin and Greek editions of the New Testament, which raised questions that would be influential in the Protestant Reformation and Catholic Counter-Reformation. He also wrote On Free Will,[7] In Praise of Folly, Handbook of a Christian Knight, On Civility in Children, Copia: Foundations of the Abundant Style, Julius Exclusus, and many other works.
+
+Erasmus lived against the backdrop of the growing European religious Reformation. While he was critical of the abuses within the Catholic Church and called for reform, he nonetheless kept his distance from Luther, Henry VIII, and John Calvin and continued to recognise the authority of the pope, emphasizing a middle way with a deep respect for traditional faith, piety and grace, and rejecting Luther's emphasis on faith alone.[citation needed] Erasmus remained a member of the Catholic Church all his life, remaining committed to reforming the Church and its clerics' abuses from within.[8][9] He also held to the doctrine of synergism, which some Reformers (Calvinists) rejected in favor of the doctrine of monergism. His middle road ("via media") approach disappointed, and even angered, scholars in both camps.
+
+Erasmus died suddenly in Basel in 1536 while preparing to return to Brabant and was buried in Basel Minster, the former cathedral of the city.[10] A bronze statue of Erasmus was erected in 1622 in his city of birth, replacing an earlier work in stone.
+
+Desiderius Erasmus is reported to have been born in Rotterdam on 28 October in the late 1460s.[3][11] He was named after Saint Erasmus of Formiae, whom Erasmus's father Gerard personally favored.[12] A 17th-century legend has it that Erasmus was first named Geert Geerts (also Gerhard Gerhards or Gerrit Gerritsz),[13] but this is unfounded.[14] A well-known wooden picture indicates: Goudæ conceptus, Roterodami natus (Latin for Conceived in Gouda, born in Rotterdam). According to an article by historian Renier Snooy (1478–1537), Erasmus was born in Gouda.
+
+The exact year of his birth is controversial but most agree it was in 1466.[15] Evidence confirming the year of Erasmus' birth in 1466 can be found in his own words: fifteen out of twenty-three statements he made about his age indicate 1466.[16] He was christened "Erasmus" after the saint of that name.[17] Although associated closely with Rotterdam, he lived there for only four years, never to return. Information on his family and early life comes mainly from vague references in his writings. His parents were not legally married. His father, Gerard, was a Catholic priest and curate in Gouda.[18] Little is known of his mother, although her known name was Margaretha Rogerius (Latinized form of Dutch surname Rutgers)[19] and she was the daughter of a doctor from Zevenbergen. She may have been Gerard's housekeeper.[15][18][20] Although he was born out of wedlock, Erasmus was cared for by his parents until their early deaths from the plague in 1483. This solidified his view of his origin as a stain and cast a pall over his youth.[18]
+
+Erasmus was given the highest education available to a young man of his day, in a series of monastic or semi-monastic schools. At the age of nine, he and his older brother Peter were sent to one of the best Latin schools in the Netherlands, located at Deventer and owned by the chapter clergy of the Lebuïnuskerk (St Lebuin's Church),[15] though some earlier biographies assert it was a school run by the Brethren of the Common Life.[15] During his stay there the curriculum was renewed by the principal of the school, Alexander Hegius. For the first time ever in Europe, Greek was taught at a lower level than a university[dubious  – discuss] and this is where he began learning it.[21] He also gleaned there the importance of a personal relationship with God but eschewed the harsh rules and strict methods of the religious brothers and educators. His education there ended when plague struck the city about 1483, and his mother, who had moved to provide a home for her sons, died from the infection.[15]
+
+Most likely in 1487,[22] poverty[23] forced Erasmus into the consecrated life as a canon regular of St. Augustine at the canonry of Stein, in South Holland. He took vows there in late 1488[22] and was ordained to the Catholic priesthood on 25 April 1492.[23] It is said that he never seemed to have actively worked as a priest for a long time,[24] and certain abuses in religious orders were among the chief objects of his later calls to reform the Church from within.
+
+While at Stein, Erasmus met with a fellow canon, Servatius Rogerus,[25] and wrote a series of passionate letters in which he called Rogerus "half my soul".[26][27] This correspondence contrasts sharply with the generally detached and much more restrained attitude he showed in his later life. Later, while tutoring in Paris, he was suddenly dismissed by the guardian of Thomas Grey. Some have taken this as evidence of an illicit affair.[28] No personal denunciation was made of Erasmus during his lifetime, however, and he took pains in later life to distance these earlier episodes by condemning sodomy in his works, and praising sexual desire in marriage between men and women.[29]
+
+Soon after his priestly ordination, he got his chance to leave the canonry when offered the post of secretary to the Bishop of Cambrai, Henry of Bergen, on account of his great skill in Latin and his reputation as a man of letters. To allow him to accept that post, he was given a temporary dispensation from his religious vows on the grounds of poor health and love of Humanistic studies, though he remained a priest. Pope Leo X later made the dispensation permanent, a considerable privilege at the time.
+
+In 1495, with Bishop Henry's consent and a stipend, Erasmus went on to study at the University of Paris in the Collège de Montaigu, a centre of reforming zeal, under the direction of the ascetic Jan Standonck, of whose rigors he complained. The University was then the chief seat of Scholastic learning but already coming under the influence of Renaissance humanism. For instance, Erasmus became an intimate friend of an Italian humanist Publio Fausto Andrelini, poet and "professor of humanity" in Paris.
+
+The chief centres of Erasmus's activity were Paris, Leuven (in the Duchy of Brabant, now in Belgium), England, and Basel; yet he never belonged firmly in any one of these places. In 1499 he was invited to England by William Blount, 4th Baron Mountjoy, who offered to accompany him on his trip to England. According to Thomas Penn, Erasmus was "ever susceptible to the charms of attractive, well-connected, and rich young men".[30] His time in England was fruitful in the making of lifelong friendships with the leaders of English thought in the days of King Henry VIII: John Colet, Thomas More, John Fisher, Thomas Linacre and William Grocyn. At the University of Cambridge, he was the Lady Margaret's Professor of Divinity and had the option of spending the rest of his life as an English professor. He stayed at Queens' College, Cambridge, from 1510 to 1515.[31] His rooms were in the "I" staircase of Old Court, and he famously hated English ale and English weather.
+
+Erasmus suffered from poor health and complained that Queens' could not supply him with enough decent wine (wine was the Renaissance medicine for gallstones, from which Erasmus suffered). Until the early 20th century, Queens' College used to have a corkscrew that was purported to be "Erasmus' corkscrew", which was a third of a metre long; even today the college still has what it calls "Erasmus' chair".[32] Today Queens' College also has an Erasmus Building and an Erasmus Room. His legacy is marked for someone who complained bitterly about the lack of comforts and luxuries to which he was accustomed. As Queens' was an unusually humanist-leaning institution in the 16th century, Queens' College Old Library still houses many first editions of Erasmus' publications, many of which were acquired during that period by bequest or purchase, including Erasmus' New Testament translation, which is signed by friend and Polish religious reformer Jan Łaski.[33] Erasmus' friend, Chancellor John Fisher, was president of Queens' College from 1505 to 1508. His friendship with Fisher is the reason he chose to stay at Queens' while lecturing in Greek at the University.[34]
+
+During his first visit to England in 1499, he taught at the University of Oxford. Erasmus was particularly impressed by the Bible teaching of John Colet, who pursued a style more akin to the church fathers than the Scholastics. This prompted him, upon his return from England, to master the Greek language, which would enable him to study theology on a more profound level and to prepare a new edition of Jerome's Bible translation. On one occasion he wrote to Colet:
+
+I cannot tell you, dear Colet, how I hurry on, with all sails set, to holy literature. How I dislike everything that keeps me back, or retards me.[23]
+
+Despite a chronic shortage of money, he succeeded in learning Greek by an intensive, day-and-night study of three years, continuously begging in letters that his friends send him books and money for teachers.[35] Discovery in 1506 of Lorenzo Valla's New Testament Notes encouraged Erasmus to continue the study of the New Testament.[36]
+
+Erasmus preferred to live the life of an independent scholar and made a conscious effort to avoid any actions or formal ties that might inhibit his freedom of intellect and literary expression. Throughout his life, he was offered positions of honor and profit in academia but declined them all, preferring the uncertain but sufficient rewards of independent literary activity. He did however assist his friend John Colet by authoring Greek textbooks and procuring members of staff for the newly established St Paul's School.[37] From 1506 to 1509, he was in Italy: in 1506 he graduated as Doctor of Divinity at the Turin University, and he spent part of the time as a proofreader at the publishing house of Aldus Manutius in Venice. According to his letters, he was associated with the Venetian natural philosopher, Giulio Camillo,[38] but, apart from this, he had a less active association with Italian scholars than might have been expected.
+
+His residence at Leuven, where he lectured at the University, exposed Erasmus to much criticism from those ascetics, academics and clerics hostile to the principles of literary and religious reform and to the loose norms of the Renaissance adherents to which he was devoting his life. In 1517, he supported the foundation at the University, by his friend Hieronymus van Busleyden, of the Collegium Trilingue for the study of Hebrew, Latin, and Greek – after the model of the College of the Three Languages at the University of Alcalá. However, feeling that the lack of sympathy that prevailed at Leuven at that time was actually a form of mental persecution, he sought refuge in Basel, where under the shelter of Swiss hospitality he could express himself freely. Admirers from all quarters of Europe visited him there and he was surrounded by devoted friends, notably developing a lasting association with the great publisher Johann Froben.
+
+Only when he had mastered Latin did he begin to express himself on major contemporary themes in literature and religion. He felt called upon to use his learning in a purification of doctrine by returning to the historic documents and original languages of sacred Scripture. He tried to free the methods of scholarship from the rigidity and formalism of medieval traditions, but he was not satisfied with this. His revolt against certain forms of Christian monasticism and scholasticism was not based on doubts about the truth of doctrine, nor from hostility to the organization of the Church itself, nor from rejection of celibacy or monastic lifestyles. He saw himself as a preacher of righteousness by an appeal to reason, applied frankly and without fear of the magisterium. He always intended to remain faithful to Catholic doctrine and therefore was convinced he could criticize frankly virtually everyone and everything. Aloof from entangling obligations, Erasmus was the centre of the literary movement of his time, corresponding with more than five hundred men in the worlds of politics and of thought.
+
+In 1502, in Spain, Cardinal Francisco Jiménez de Cisneros had put together a team of Spanish translators to create a Compilation of a Bible in four languages, Greek, Hebrew, Aramaic and Latin. Translators for Greek were commissioned from Greece itself and worked closely with Latinists. Cardinal Cisneros's team completed and printed the full New Testament, including the Greek translation, in 1514. To do so they developed specific types to print Greek. Cisneros informed Erasmus of the works going on in Spain and may have sent a printed version of the New Testament to him. However, the Spanish team wanted the entire Bible to be released as one single work and withdrew from publication.
+
+The information and the delay allowed Erasmus to request a "Publication Privilege" of four years for the Greek New Testament to ensure that his work would be published first. He obtained it in 1516 from both Pope Leo X, to whom he would dedicate his work, and Emperor Maximilian I. Erasmus' Greek New Testament was published first, in 1516, forcing the Spanish team of Cisneros to wait until 1520 to publish their Complutensian Polyglot Bible.[39][40]
+
+It is hard to say if Erasmus' actions had an effect on delaying the publication of Complutensian Polyglot, causing the Spanish team to take more time, or if it made no difference in their perfectionism. The Spanish copy was approved for publication by the Pope in 1520; however, it was not released until 1522 due to the team's insistence on reviewing and editing. Only fifteen errors have been found in the entire six volumes and four languages of Cisneros's bible, an extraordinarily low number for the time. The fear of them publishing first, though, affected Erasmus's work, rushing him to printing and causing him to forego editing. The result was a large number of translation mistakes, transcription errors, and typos, that required further editions to be printed (see "publication").[41]
+
+Erasmus had been working for years on two projects: a collation of Greek texts and a fresh Latin New Testament. In 1512, he began his work on this Latin New Testament. He collected all the Vulgate manuscripts he could find to create a critical edition. Then he polished the language. He declared, "It is only fair that Paul should address the Romans in somewhat better Latin."[42] In the earlier phases of the project, he never mentioned a Greek text:
+
+My mind is so excited at the thought of emending Jerome’s text, with notes, that I seem to myself inspired by some god. I have already almost finished emending him by collating a large number of ancient manuscripts, and this I am doing at enormous personal expense.[43]
+
+While his intentions for publishing a fresh Latin translation are clear, it is less clear why he included the Greek text. Though some speculate that he intended to produce a critical Greek text or that he wanted to beat the Complutensian Polyglot into print, there is no evidence to support this. He wrote, "There remains the New Testament translated by me, with the Greek facing, and notes on it by me."[44] He further demonstrated the reason for the inclusion of the Greek text when defending his work:
+
+But one thing the facts cry out, and it can be clear, as they say, even to a blind man, that often through the translator’s clumsiness or inattention the Greek has been wrongly rendered; often the true and genuine reading has been corrupted by ignorant scribes, which we see happen every day, or altered by scribes who are half-taught and half-asleep.[45]
+
+So he included the Greek text to permit qualified readers to verify the quality of his Latin version. But by first calling the final product Novum Instrumentum omne ("All of the New Teaching") and later Novum Testamentum omne ("All of the New Testament") he also indicated clearly that he considered a text in which the Greek and the Latin versions were consistently comparable to be the essential core of the church's New Testament tradition.
+
+In a way it is legitimate to say that Erasmus "synchronized" or "unified" the Greek and the Latin traditions of the New Testament by producing an updated translation of both simultaneously. Both being part of canonical tradition, he clearly found it necessary to ensure that both were actually present in the same content. In modern terminology, he made the two traditions "compatible". This is clearly evidenced by the fact that his Greek text is not just the basis for his Latin translation, but also the other way round: there are numerous instances where he edits the Greek text to reflect his Latin version. For instance, since the last six verses of Revelation were missing from his Greek manuscript, Erasmus translated the Vulgate's text back into Greek. Erasmus also translated the Latin text into Greek wherever he found that the Greek text and the accompanying commentaries were mixed up, or where he simply preferred the Vulgate's reading to the Greek text.[46]
+
+Erasmus said it was "rushed into print rather than edited",[41] resulting in a number of transcription errors. After comparing what writings he could find, Erasmus wrote corrections between the lines of the manuscripts he was using (among which was Minuscule 2) and sent them as proofs to Froben.[47] His hurried effort was published by his friend Johann Froben of Basel in 1516 and thence became the first published Greek New Testament, the Novum Instrumentum omne, diligenter ab Erasmo Rot. Recognitum et Emendatum. Erasmus used several Greek manuscript sources because he did not have access to a single complete manuscript. Most of the manuscripts were, however, late Greek manuscripts of the Byzantine textual family and Erasmus used the oldest manuscript the least because "he was afraid of its supposedly erratic text."[48] He also ignored much older and better manuscripts that were at his disposal.[49]
+
+In the second (1519) edition, the more familiar term Testamentum was used instead of Instrumentum. This edition was used by Martin Luther in his German translation of the Bible, written for people who could not understand Latin. Together, the first and second editions sold 3,300 copies. By comparison, only 600 copies of the Complutensian Polyglot were ever printed. The first and second edition texts did not include the passage (1 John 5:7–8) that has become known as the Comma Johanneum. Erasmus had been unable to find those verses in any Greek manuscript, but one was supplied to him during production of the third edition.  The Catholic Church decreed that the Comma Johanneum was open to dispute (2 June 1927), and it is rarely included in modern scholarly translations.
+
+The third edition of 1522 was probably used by Tyndale for the first English New Testament (Worms, 1526) and was the basis for the 1550 Robert Stephanus edition used by the translators of the Geneva Bible and King James Version of the English Bible. Erasmus published a fourth edition in 1527 containing parallel columns of Greek, Latin Vulgate and Erasmus' Latin texts. In this edition Erasmus also supplied the Greek text of the last six verses of Revelation (which he had translated from Latin back into Greek in his first edition) from Cardinal Ximenez's Biblia Complutensis. In 1535 Erasmus published the fifth (and final) edition which dropped the Latin Vulgate column but was otherwise similar to the fourth edition. Later versions of the Greek New Testament by others, but based on Erasmus' Greek New Testament, became known as the Textus Receptus.[50]
+
+Erasmus dedicated his work to Pope Leo X as a patron of learning and regarded this work as his chief service to the cause of Christianity. Immediately afterwards, he began the publication of his Paraphrases of the New Testament, a popular presentation of the contents of the several books. These, like all of his writings, were published in Latin but were quickly translated into other languages with his encouragement.
+
+Erasmus, in his capacity as humanist editor, advised major printers such as Aldus Manutius on which manuscripts to publish.[51]
+
+The Protestant Reformation began in the year following the publication of his edition of the Greek New Testament (1516) and tested Erasmus' character. The issues between the Catholic Church and the growing religious movement which would later become known as Protestantism, had become so clear that few could escape the summons to join the debate. Erasmus, at the height of his literary fame, was inevitably called upon to take sides, but partisanship was foreign to his nature and his habits. Despite all his criticism of clerical corruption and abuses within the Catholic Church,[8] which lasted for years and was also directed towards many of the Church's basic doctrines,[9] Erasmus shunned the Reformation movement along with its most radical and reactionary offshoots,[8][need quotation to verify] and sided with neither party.[8]
+
+The world had laughed at his satire, but few had interfered with his activities. He believed that his work so far had commended itself to the best minds and also to the dominant powers in the religious world. Erasmus did not build a large body of supporters with his letters. He chose to write in Greek and Latin, the languages of scholars. His critiques reached an elite but small audience.[52]
+
+"Free will does not exist", according to Luther in his letter De Servo Arbitrio to Erasmus translated into German by Justus Jonas (1526), in that sin makes human beings completely incapable of bringing themselves to God. Noting Luther's criticism of the Catholic Church, Erasmus described him as "a mighty trumpet of gospel truth" while agreeing, "It is clear that many of the reforms for which Luther calls are urgently needed."[53] He had great respect for Luther, and Luther spoke with admiration of Erasmus' superior learning. Luther hoped for his cooperation in a work which seemed only the natural outcome of his own. In their early correspondence, Luther expressed boundless admiration for all Erasmus had done in the cause of a sound and reasonable Christianity and urged him to join the Lutheran party. Erasmus declined to commit himself, arguing that to do so would endanger his position as a leader in the movement for pure scholarship which he regarded as his purpose in life. Only as an independent scholar could he hope to influence the reform of religion. When Erasmus hesitated to support him, the straightforward Luther became angered that Erasmus was avoiding the responsibility due either to cowardice or a lack of purpose. However, any hesitancy on the part of Erasmus stemmed, not from lack of courage or conviction, but rather from a concern over the mounting disorder and violence of the reform movement. To Philip Melanchthon in 1524 he wrote:
+
+I know nothing of your church; at the very least it contains people who will, I fear, overturn the whole system and drive the princes into using force to restrain good men and bad alike. The gospel, the word of God, faith, Christ, and Holy Spirit – these words are always on their lips; look at their lives and they speak quite another language.[54]
+
+Again, in 1529, he writes “An epistle against those who falsely boast they are Evangelicals”[55] to Vulturius Neocomus (Gerardus Geldenhouwer). Here Erasmus complains of the doctrines and morals of the Reformers:
+
+You declaim bitterly against the luxury of priests, the ambition of bishops, the tyranny of the Roman Pontiff, and the babbling of the sophists; against our prayers, fasts, and Masses; and you are not content to retrench the abuses that may be in these things, but must needs abolish them entirely. ...Look around on this ‘Evangelical’ generation,[56] and observe whether amongst them less indulgence is given to luxury, lust, or avarice, than amongst those whom you so detest. Show me any one person who by that Gospel has been reclaimed from drunkenness to sobriety, from fury and passion to meekness, from avarice to liberality, from reviling to well-speaking, from wantonness to modesty. I will show you a great many who have become worse through following it. ...The solemn prayers of the Church are abolished, but now there are very many who never pray at all. ...I have never entered their conventicles, but I have sometimes seen them returning from their sermons, the countenances of all of them displaying rage, and wonderful ferocity, as though they were animated by the evil spirit. ...Who ever beheld in their meetings any one of them shedding tears, smiting his breast, or grieving for his sins? ...Confession to the priest is abolished, but very few now confess to God. ...They have fled from Judaism that they may become Epicureans.[57]
+
+Apart from these perceived moral failings of the Reformers, Erasmus also dreaded any change in doctrine, citing the long history of the Church as a bulwark against innovation. In book I of his Hyperaspistes he puts the matter bluntly to Luther:
+
+We are dealing with this: Would a stable mind depart from the opinion handed down by so many men famous for holiness and miracles, depart from the decisions of the Church, and commit our souls to the faith of someone like you who has sprung up just now with a few followers, although the leading men of your flock do not agree either with you or among themselves – indeed though you do not even agree with yourself, since in this same Assertion[58] you say one thing in the beginning and something else later on, recanting what you said before.[59]
+
+Continuing his chastisement of Luther – and undoubtedly put off by the notion of there being "no pure interpretation of Scripture anywhere but in Wittenberg"[60] – Erasmus touches upon another important point of the controversy:
+
+You stipulate that we should not ask for or accept anything but Holy Scripture, but you do it in such a way as to require that we permit you to be its sole interpreter, renouncing all others. Thus the victory will be yours if we allow you to be not the steward but the lord of Holy Scripture.[61]
+
+Though he sought to remain firmly neutral in doctrinal disputes, each side accused him of siding with the other, perhaps because of his neutrality. It was not for lack of fidelity with either side but a desire for fidelity with them both:
+
+I detest dissension because it goes both against the teachings of Christ and against a secret inclination of nature. I doubt that either side in the dispute can be suppressed without grave loss.[53]
+
+In his catechism (entitled Explanation of the Apostles' Creed) (1533), Erasmus took a stand against Luther's teaching by asserting the unwritten Sacred Tradition as just as valid a source of revelation as the Bible, by enumerating the Deuterocanonical books in the canon of the Bible and by acknowledging seven sacraments.[62] He called "blasphemers" anyone who questioned the perpetual virginity of Mary.[63] However, he supported lay access to the Bible.[63]
+
+In a letter to Nikolaus von Amsdorf, Luther objected to Erasmus’ catechism and called Erasmus a "viper," "liar," and "the very mouth and organ of Satan".[64]
+
+As regards the Reformation, Erasmus was accused by the monks to have:
+
+prepared the way and was responsible for Martin Luther. Erasmus, they said, had laid the egg, and Luther had hatched it. Erasmus wittily dismissed the charge, claiming that Luther had hatched a different bird entirely.[65]
+
+Twice in the course of the great discussion, he allowed himself to enter the field of doctrinal controversy, a field foreign to both his nature and his previous practice. One of the topics he dealt with was free will, a crucial question. In his De libero arbitrio diatribe sive collatio (1524), he lampoons the Lutheran view on free will. He lays down both sides of the argument impartially. The "Diatribe" did not encourage any definite action; this was its merit to the Erasmians and its fault in the eyes of the Lutherans. In response, Luther wrote his De servo arbitrio (On the Bondage of the Will, 1525), which attacks the "Diatribe" and Erasmus himself, going so far as to claim that Erasmus was not a Christian. Erasmus responded with a lengthy, two-part Hyperaspistes (1526–27). In this controversy Erasmus lets it be seen that he would like to claim more for free will than St. Paul and St. Augustine seem to allow according to Luther's interpretation.[66] For Erasmus the essential point is that humans have the freedom of choice.[67] The conclusions Erasmus reached drew upon a large array of notable authorities, including, from the Patristic period, Origen, John Chrysostom, Ambrose, Jerome, and Augustine, in addition to many leading Scholastic authors, such as Thomas Aquinas and Duns Scotus. The content of Erasmus' works also engaged with later thought on the state of the question, including the perspectives of the via moderna school and of Lorenzo Valla, whose ideas he rejected.
+
+As the popular response to Luther gathered momentum, the social disorders, which Erasmus dreaded and Luther disassociated himself from, began to appear, including the German Peasants' War, the Anabaptist disturbances in Germany and in the Low Countries, iconoclasm, and the radicalization of peasants across Europe. If these were the outcomes of reform, he was thankful that he had kept out of it. Yet he was ever more bitterly accused of having started the whole "tragedy" (as the Catholics dubbed Protestantism).
+
+When the city of Basel definitely adopted the Reformation in 1529, Erasmus gave up his residence there and settled in the imperial town of Freiburg im Breisgau.
+
+Certain works of Erasmus laid a foundation for religious toleration and ecumenism. For example, in De libero arbitrio, opposing certain views of Martin Luther, Erasmus noted that religious disputants should be temperate in their language, "because in this way the truth, which is often lost amidst too much wrangling may be more surely perceived." Gary Remer writes, "Like Cicero, Erasmus concludes that truth is furthered by a more harmonious relationship between interlocutors."[68] Although Erasmus did not oppose the punishment of heretics, in individual cases he generally argued for moderation and against the death penalty. He wrote, "It is better to cure a sick man than to kill him."[69]
+
+A test of the Reformation was the doctrine of the sacraments, and the crux of this question was the observance of the Eucharist. In 1530, Erasmus published a new edition of the orthodox treatise of Algerus against the heretic Berengar of Tours in the eleventh century. He added a dedication, affirming his belief in the reality of the Body of Christ after consecration in the Eucharist, commonly referred to as transubstantiation. The sacramentarians, headed by Œcolampadius of Basel, were, as Erasmus says, quoting him as holding views similar to their own in order to try to claim him for their schismatic and "erroneous" movement.
+
+When his strength began to fail, he decided to accept an invitation by Queen Mary of Hungary, Regent of the Netherlands, to move from Freiburg to Brabant. However, during preparations for the move in 1536, he suddenly died from an attack of dysentery during a visit to Basel.[70] He had remained loyal to the papal authorities in Rome, but he did not have the opportunity to receive the last rites of the Catholic Church; the reports of his death do not mention whether he asked for a priest or not. According to Jan van Herwaarden, this is consistent with his view that outward signs were not important; what mattered is the believer's direct relationship with God, which he noted "as the [Catholic] church believes". However, Herwaarden observes that "he did not dismiss the rites and sacraments out of hand but asserted a dying person could achieve a state of salvation without the priestly rites, provided their faith and spirit were attuned to God."[71]
+He was buried with great ceremony in Basel Minster (the former cathedral) there.[10]
+
+His last words, as recorded by his friend Beatus Rhenanus, were apparently "Dear God" (Dutch: Lieve God).[72] A bronze statue of him was erected in the city of his birth in 1622, replacing an earlier work in stone.
+
+Erasmus wrote both on church subjects and those of general human interest. By the 1530s, the writings of Erasmus accounted for 10 to 20 percent of all book sales in Europe.[73]  He is credited with coining the adage, "In the land of the blind, the one-eyed man is king."  With the collaboration of Publio Fausto Andrelini, he formed a paremiography (collection) of Latin proverbs and adages, commonly titled Adagia.  Erasmus is also generally credited with originating the phrase "Pandora's box", arising through an error in his translation of Hesiod's Pandora in which he confused pithos (storage jar) with pyxis (box).
+
+His more serious writings begin early with the Enchiridion militis Christiani, the "Handbook of the Christian Soldier" (1503 – translated into English a few years later by the young William Tyndale).  (A more literal translation of enchiridion – "dagger" – has been likened to "the spiritual equivalent of the modern Swiss Army knife.")[74] In this short work, Erasmus outlines the views of the normal Christian life, which he was to spend the rest of his days elaborating.  The chief evil of the day, he says, is formalism – going through the motions of tradition without understanding their basis in the teachings of Christ.  Forms can teach the soul how to worship God, or they may hide or quench the spirit. In his examination of the dangers of formalism, Erasmus discusses monasticism, saint worship, war, the spirit of class and the foibles of "society."
+
+The Enchiridion is more like a sermon than a satire. With it Erasmus challenged common assumptions, painting the clergy as educators who should share the treasury of their knowledge with the laity. He emphasized personal spiritual disciplines and called for a reformation which he characterized as a collective return to the Fathers and Scripture. Most importantly, he extolled the reading of scripture as vital because of its power to transform and motivate toward love. Much like the Brethren of the Common Life, he wrote that the New Testament is the law of Christ people are called to obey and that Christ is the example they are called to imitate.
+
+According to Ernest Barker, "Besides his work on the New Testament, Erasmus laboured also, and even more arduously, on the early Fathers. …Among the Latin Fathers he edited the works of St Jerome, St Hilary, and St Augustine; among the Greeks he worked on Irenaeus, Origen and Chrysostom."[75]
+
+Erasmus also wrote of the legendary Frisian freedom fighter and rebel Pier Gerlofs Donia (Greate Pier), though more often in criticism than in praise of his exploits.  Erasmus saw him as a dim, brutal man who preferred physical strength to wisdom.[76]
+
+One of Erasmus' best-known works, inspired by De triumpho stultitiae (written by Italian humanist Faustino Perisauli), is The Praise of Folly, published under the double title Moriae encomium (Greek, Latinised) and Laus stultitiae (Latin).[77] A satirical attack on superstitions and other traditions of European society in general and the western Church in particular, it was written in 1509, published in 1511, and dedicated to Sir Thomas More, whose name the title puns.
+
+The Institutio principis Christiani  or "Education of a Christian Prince" (Basel, 1516) was written as advice to the young king Charles of Spain (later Charles V, Holy Roman Emperor). Erasmus applies the general principles of honor and sincerity to the special functions of the Prince, whom he represents throughout as the servant of the people. Education was published in 1516, three years after[78] Niccolò Machiavelli’s The Prince was written; a comparison between the two is worth noting. Machiavelli stated that, to maintain control by political force, it is safer for a prince to be feared than loved. Erasmus preferred for the prince to be loved, and strongly suggested a well-rounded education in order to govern justly and benevolently and avoid becoming a source of oppression.
+
+As a result of his reformatory activities, Erasmus found himself at odds with both of the great parties. His last years were embittered by controversies with men toward whom he was sympathetic. Notable among these was Ulrich von Hutten, a brilliant but erratic genius who had thrown himself into the Lutheran cause and declared that Erasmus, if he had a spark of honesty, would do the same. In his reply in 1523, Spongia adversus aspergines Hutteni, Erasmus displays his skill in semantics. He accuses Hutten of having misinterpreted his utterances about reform and reiterates his determination never to break with the Church.
+
+The Ciceronianus came out in 1528, attacking the style of Latin that was based exclusively and fanatically on Cicero's writings. Etienne Dolet wrote a riposte titled Erasmianus in 1535.
+
+Erasmus' last major work, published the year of his death, is the Ecclesiastes or "Gospel Preacher" (Basel, 1536), a massive manual for preachers of around a thousand pages. Though somewhat unwieldy because Erasmus was unable to edit it properly in his old age, it is in some ways the culmination of all of Erasmus' literary and theological learning, offering prospective preachers advice on nearly every conceivable aspect of their vocation with extraordinarily abundant reference to classical and biblical sources.
+
+Erasmus’ Sileni Alcibiadis is one of his most direct assessments of the need for Church reform. Johann Froben published it first within a revised edition of the Adagia in 1515, then as a stand-alone work in 1517. This essay has been likened to John Colet's Convocation Sermon, though the styles differ.
+
+Sileni is the plural (Latin) form of Silenus, a creature often related to the Roman wine god Bacchus and represented in pictorial art as inebriated, merry revellers, variously mounted on donkeys, singing, dancing, playing flutes, etc. Alcibiades was a Greek politician in the 5th century BCE and a general in the Peloponnesian War; he figures here more as a character written into some of Plato's dialogues – a young, debauched playboy whom Socrates tries to convince to seek truth instead of pleasure, wisdom instead of pomp and splendor.
+
+The term Sileni – especially when juxtaposed with the character of Alcibiades – can therefore be understood as an evocation of the notion that something on the inside is more expressive of a person's character than what one sees on the outside. For instance, something or someone ugly on the outside can be beautiful on the inside, which is one of the main points of Plato's dialogues featuring Alcibiades and the Symposion, in which Alcibiades also appears.
+
+In support of this, Erasmus states, "Anyone who looks closely at the inward nature and essence will find that nobody is further from true wisdom than those people with their grand titles, learned bonnets, splendid sashes and bejeweled rings, who profess to be wisdom’s peak."  On the other hand, Erasmus lists several Sileni and then questions whether Christ is the most noticeable Silenus of them all. The Apostles were Sileni since they were ridiculed by others. He believes that the things which are the least ostentatious can be the most significant, and that the Church constitutes all Christian people – that despite contemporary references to clergy as the whole of the Church, they are merely its servants. He criticizes those that spend the Church's riches at the people's expense. The true point of the Church is to help people lead Christian lives. Priests are supposed to be pure, yet when they stray, no one condemns them. He criticizes the riches of the popes, believing that it would be better for the Gospel to be most important.
+
+The popularity of his books is reflected in the number of editions and translations that have appeared since the sixteenth century. Ten columns of the catalogue of the British Library are taken up with the enumeration of the works and their subsequent reprints. The greatest names of the classical and patristic world are among those translated, edited, or annotated by Erasmus, including Saint Ambrose, Aristotle, Saint Augustine, Saint Basil, Saint John Chrysostom, Cicero and Saint Jerome.[79]
+
+In his native Rotterdam, the University and Gymnasium Erasmianum have been named in his honor. In 2003, a poll showing that most Rotterdammers believed Erasmus to be the designer of the local "Erasmus Bridge" instigated the founding of the Erasmus House (Rotterdam),[80] and the Erasmus House (Jakarta)[81] dedicated to celebrating Erasmus' legacy. Three moments in Erasmus' life are celebrated annually. On 1 April, the city celebrates the publication of his best-known book The Praise of Folly. On 11 July, the Night of Erasmus celebrates the lasting influence of his work. His birthday is celebrated on 28 October.[82]
+
+Erasmus' reputation and the interpretations of his work have varied over time. Moderate Catholics recognized him as a leading figure in attempts to reform the Church, while Protestants recognized his initial support for Luther's ideas and the groundwork he laid for the future Reformation, especially in biblical scholarship. By the 1560s, however, there was a marked change in reception.[citation needed]
+
+According to Franz Anton Knittel, Erasmus in his Novum Instrumentum omne did not incorporate the Comma from the Codex Montfortianus (concerning the Trinity), because of grammar differences, but used Complutensian Polyglotta. According to him the Comma was known to Tertullian.[83]
+
+Protestant views of Erasmus fluctuated depending on region and period, with continual support in his native Netherlands and in cities of the Upper Rhine area. However, following his death and in the late sixteenth century, many Reformation supporters saw Erasmus's critiques of Luther and lifelong support for the universal Catholic Church as damning, and second-generation Protestants were less vocal in their debts to the great humanist. Nevertheless, his reception is demonstrable among Swiss Protestants in the sixteenth century: he had an indelible influence on the biblical commentaries of, for example, Konrad Pellikan, Heinrich Bullinger, and John Calvin, all of whom used both his annotations on the New Testament and his paraphrases of same in their own New Testament commentaries.[84]
+
+However, Erasmus designated his own legacy, and his life works were turned over at his death to his friend the Protestant humanist turned remonstrator Sebastian Castellio for the repair of the breach and divide of Christianity in its Catholic, Anabaptist, and Protestant branches.[85]
+
+By the coming of the Age of Enlightenment, however, Erasmus increasingly again became a more widely respected cultural symbol and was hailed as an important figure by increasingly broad groups. In a letter to a friend, Erasmus once had written: "That you are patriotic will be praised by many and easily forgiven by everyone; but in my opinion it is wiser to treat men and things as though we held this world the common fatherland of all."[86] Thus, the universalist ideals of Erasmus are sometimes claimed to be important for fixing global governance.[87]
+
+Several schools, faculties and universities in the Netherlands and Belgium are named after him, as is Erasmus Hall in Brooklyn, New York, USA.The European Union's Erasmus Programme scholarships enable students to spend up to a year of their university courses in a university in another European country.
+
+Eramus is credited with saying "When I get a little money I buy books; and if any is left, I buy food and clothes."[88]
+
+He is also blamed for the mistranslation from Greek of to call a bowl a bowl as to call a spade a spade.[89]
+

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+A desert is a barren area of landscape where little precipitation occurs and, consequently, living conditions are hostile for plant and animal life. The lack of vegetation exposes the unprotected surface of the ground to the processes of denudation. About one-third of the land surface of the world is arid or semi-arid. This includes much of the polar regions, where little precipitation occurs, and which are sometimes called polar deserts or "cold deserts". Deserts can be classified by the amount of precipitation that falls, by the temperature that prevails, by the causes of desertification or by their geographical location.
+
+Deserts are formed by weathering processes as large variations in temperature between day and night put strains on the rocks, which consequently break in pieces. Although rain seldom occurs in deserts, there are occasional downpours that can result in flash floods. Rain falling on hot rocks can cause them to shatter, and the resulting fragments and rubble strewn over the desert floor are further eroded by the wind. This picks up particles of sand and dust and wafts them aloft in sand or dust storms. Wind-blown sand grains striking any solid object in their path can abrade the surface. Rocks are smoothed down, and the wind sorts sand into uniform deposits. The grains end up as level sheets of sand or are piled high in billowing sand dunes. Other deserts are flat, stony plains where all the fine material has been blown away and the surface consists of a mosaic of smooth stones. These areas are known as desert pavements, and little further erosion takes place. Other desert features include rock outcrops, exposed bedrock and clays once deposited by flowing water. Temporary lakes may form and salt pans may be left when waters evaporate. There may be underground sources of water, in the form of springs and seepages from aquifers. Where these are found, oases can occur.
+
+Plants and animals living in the desert need special adaptations to survive in the harsh environment. Plants tend to be tough and wiry with small or no leaves, water-resistant cuticles, and often spines to deter herbivory. Some annual plants germinate, bloom and die in the course of a few weeks after rainfall, while other long-lived plants survive for years and have deep root systems able to tap underground moisture. Animals need to keep cool and find enough food and water to survive. Many are nocturnal, and stay in the shade or underground during the heat of the day. They tend to be efficient at conserving water, extracting most of their needs from their food and concentrating their urine. Some animals remain in a state of dormancy for long periods, ready to become active again during the rare rainfall. They then reproduce rapidly while conditions are favorable before returning to dormancy.
+
+People have struggled to live in deserts and the surrounding semi-arid lands for millennia. Nomads have moved their flocks and herds to wherever grazing is available, and oases have provided opportunities for a more settled way of life. The cultivation of semi-arid regions encourages erosion of soil and is one of the causes of increased desertification. Desert farming is possible with the aid of irrigation, and the Imperial Valley in California provides an example of how previously barren land can be made productive by the import of water from an outside source. Many trade routes have been forged across deserts, especially across the Sahara Desert, and traditionally were used by caravans of camels carrying salt, gold, ivory and other goods. Large numbers of slaves were also taken northwards across the Sahara[citation needed]. Some mineral extraction also takes place in deserts, and the uninterrupted sunlight gives potential for the capture of large quantities of solar energy.
+
+English desert and its Romance cognates (including Italian and Portuguese deserto, French désert and Spanish desierto) all come from the ecclesiastical Latin dēsertum (originally "an abandoned place"), a participle of dēserere, "to abandon".[1] The correlation between aridity and sparse population is complex and dynamic, varying by culture, era, and technologies; thus the use of the word desert can cause confusion. In English before the 20th century, desert was often used in the sense of "unpopulated area", without specific reference to aridity;[1] but today the word is most often used in its climate-science sense (an area of low precipitation).[2] Phrases such as "desert island"[3] and "Great American Desert", or Shakespeare's "deserts of Bohemia" (The Winter's Tale) in previous centuries did not necessarily imply sand or aridity; their focus was the sparse population.[4]
+
+A desert is a region of land that is very dry because it receives low amounts of precipitation (usually in the form of rain, but it may be snow, mist or fog), often has little coverage by plants, and in which streams dry up unless they are supplied by water from outside the area.[5] Deserts generally receive less than 250 mm (10 in) of precipitation each year.[5] The potential evapotranspiration may be large but (in the absence of available water) the actual evapotranspiration may be close to zero.[6] Semi-deserts are regions which receive between 250 and 500 mm (10 and 20 in) and when clad in grass, these are known as steppes.[7][8]
+
+Deserts have been defined and classified in a number of ways, generally combining total precipitation, number of days on which this falls, temperature, and humidity, and sometimes additional factors.[8] For example, Phoenix, Arizona, receives less than 250 mm (9.8 in) of precipitation per year, and is immediately recognized as being located in a desert because of its aridity-adapted plants. The North Slope of Alaska's Brooks Range also receives less than 250 mm (9.8 in) of precipitation per year and is often classified as a cold desert.[9] Other regions of the world have cold deserts, including areas of the Himalayas[10] and other high-altitude areas in other parts of the world.[11] Polar deserts cover much of the ice-free areas of the Arctic and Antarctic.[12][13] A non-technical definition is that deserts are those parts of the Earth's surface that have insufficient vegetation cover to support a human population.[14]
+
+Potential evapotranspiration supplements the measurement of precipitation in providing a scientific measurement-based definition of a desert. The water budget of an area can be calculated using the formula P − PE ± S, wherein P is precipitation, PE is potential evapotranspiration rates and S is the amount of surface storage of water. Evapotranspiration is the combination of water loss through atmospheric evaporation and through the life processes of plants. Potential evapotranspiration, then, is the amount of water that could evaporate in any given region. As an example, Tucson, Arizona receives about 300 mm (12 in) of rain per year, however about 2,500 mm (98 in) of water could evaporate over the course of a year.[15] In other words, about eight times more water could evaporate from the region than actually falls as rain. Rates of evapotranspiration in cold regions such as Alaska are much lower because of the lack of heat to aid in the evaporation process.[16]
+
+Deserts are sometimes classified as "hot" or "cold", "semiarid" or "coastal".[14] The characteristics of hot deserts include high temperatures in summer; greater evaporation than precipitation, usually exacerbated by high temperatures, strong winds and lack of cloud cover; considerable variation in the occurrence of precipitation, its intensity and distribution; and low humidity. Winter temperatures vary considerably between different deserts and are often related to the location of the desert on the continental landmass and the latitude. Daily variations in temperature can be as great as 22 °C (40 °F) or more, with heat loss by radiation at night being increased by the clear skies.[17]
+
+Cold deserts, sometimes known as temperate deserts, occur at higher latitudes than hot deserts, and the aridity is caused by the dryness of the air. Some cold deserts are far from the ocean and others are separated by mountain ranges from the sea, and in both cases, there is insufficient moisture in the air to cause much precipitation. The largest of these deserts are found in Central Asia. Others occur on the eastern side of the Rocky Mountains, the eastern side of the southern Andes and in southern Australia.[7] Polar deserts are a particular class of cold desert. The air is very cold and carries little moisture so little precipitation occurs and what does fall, usually snow, is carried along in the often strong wind and may form blizzards, drifts and dunes similar to those caused by dust and sand in other desert regions. In Antarctica, for example, the annual precipitation is about 50 mm (2 in) on the central plateau and some ten times that amount on some major peninsulas.[17]
+
+Based on precipitation alone, hyperarid deserts receive less than 25 mm (1 in) of rainfall a year; they have no annual seasonal cycle of precipitation and experience twelve-month periods with no rainfall at all.[17][18] Arid deserts receive between 25 and 200 mm (1 and 8 in) in a year and semiarid deserts between 200 and 500 mm (8 and 20 in). However, such factors as the temperature, humidity, rate of evaporation and evapotranspiration, and the moisture storage capacity of the ground have a marked effect on the degree of aridity and the plant and animal life that can be sustained. Rain falling in the cold season may be more effective at promoting plant growth, and defining the boundaries of deserts and the semiarid regions that surround them on the grounds of precipitation alone is problematic.[17]
+
+A semi-arid desert or a steppe is a version of the arid desert with much more rainfall, vegetation and higher humidity. These regions feature a semi-arid climate and are less extreme than regular deserts.[19] Like arid deserts, temperatures can vary greatly in semi deserts. They share some characteristics of a true desert and are usually located at the edge of deserts and continental dry areas. They usually receive precipitation from 250 mm (10 in) to 500 mm (20 in) but this can vary due to evapotranspiration and soil nutrition. Semi deserts can be found in the Tabernas Desert (and some of the Spanish Plateau), The Sahel, The Eurasian Steppe, most of Central Asia, the Western US, most of Northern Mexico, portions of South America (especially in Argentina) and the Australian Outback.[20] They usually feature BSh (hot steppe) or BSk (temperate steppe) in the Köppen climate classification.
+
+Coastal deserts are mostly found on the western edges of continental land masses in regions where cold currents approach the land or cold water upwellings rise from the ocean depths. The cool winds crossing this water pick up little moisture and the coastal regions have low temperatures and very low rainfall, the main precipitation being in the form of fog and dew. The range of temperatures on a daily and annual scale is relatively low, being 11 °C (20 °F) and 5 °C (9 °F) respectively in the Atacama Desert. Deserts of this type are often long and narrow and bounded to the east by mountain ranges. They occur in Namibia, Chile, southern California and Baja California. Other coastal deserts influenced by cold currents are found in Western Australia, the Arabian Peninsula and Horn of Africa, and the western fringes of the Sahara.[17]
+
+In 1961, Peveril Meigs divided desert regions on Earth into three categories according to the amount of precipitation they received. In this now widely accepted system, extremely arid lands have at least twelve consecutive months without precipitation, arid lands have less than 250 mm (10 in) of annual precipitation, and semiarid lands have a mean annual precipitation of between 250 and 500 mm (10–20 in). Both extremely arid and arid lands are considered to be deserts while semiarid lands are generally referred to as steppes when they are grasslands.[8]
+
+Deserts are also classified, according to their geographical location and dominant weather pattern, as trade wind, mid-latitude, rain shadow, coastal, monsoon, or polar deserts.[21] Trade wind deserts occur either side of the horse latitudes at 30° to 35° North and South. These belts are associated with the subtropical anticyclone and the large-scale descent of dry air moving from high-altitudes toward the poles. The Sahara Desert is of this type.[22] Mid-latitude deserts occur between 30° and 50° North and South. They are mostly in areas remote from the sea where most of the moisture has already precipitated from the prevailing winds. They include the Tengger and Sonoran Deserts.[21] Monsoon deserts are similar. They occur in regions where large temperature differences occur between sea and land. Moist warm air rises over the land, deposits its water content and circulates back to sea. Further inland, areas receive very little precipitation. The Thar Desert near the India/Pakistan border is of this type.[21]
+
+In some parts of the world, deserts are created by a rain shadow effect. Orographic lift occurs as air masses rise to pass over high ground. In the process they cool and lose much of their moisture by precipitation on the windward slope of the mountain range. When they descend on the leeward side, they warm and their capacity to hold moisture increases so an area with relatively little precipitation occurs.[23] The Taklamakan Desert is an example, lying in the rain shadow of the Himalayas and receiving less than 38 mm (1.5 in) precipitation annually.[24]
+Other areas are arid by virtue of being a very long way from the nearest available sources of moisture.[25]
+
+Montane deserts are arid places with a very high altitude; the most prominent example is found north of the Himalayas, in the Kunlun Mountains and the Tibetan Plateau. Many locations within this category have elevations exceeding 3,000 m (9,800 ft) and the thermal regime can be hemiboreal. These places owe their profound aridity (the average annual precipitation is often less than 40 mm or 1.5 in) to being very far from the nearest available sources of moisture and are often in the lee of mountain ranges. Montane deserts are normally cold, or may be scorchingly hot by day and very cold by night as is true of the northeastern slopes of Mount Kilimanjaro.[26]
+
+Polar deserts such as McMurdo Dry Valleys remain ice-free because of the dry katabatic winds that flow downhill from the surrounding mountains.[27] Former desert areas presently in non-arid environments, such as the Sandhills in Nebraska, are known as paleodeserts.[21] In the Köppen climate classification system, deserts are classed as BWh (hot desert) or BWk (temperate desert). In the Thornthwaite climate classification system, deserts would be classified as arid megathermal climates.[28][29]
+
+Deserts usually have a large diurnal and seasonal temperature range, with high daytime temperatures falling sharply at night. The diurnal range may be as much as 20 to 30 °C (36 to 54 °F) and the rock surface experiences even greater temperature differentials.[30] During the day the sky is usually clear and most of the sun's radiation reaches the ground, but as soon as the sun sets, the desert cools quickly by radiating heat into space. In hot deserts, the temperature during daytime can exceed 45 °C (113 °F) in summer and plunge below freezing point at night during winter.[31]
+
+Such large temperature variations have a destructive effect on the exposed rocky surfaces. The repeated fluctuations put a strain on exposed rock and the flanks of mountains crack and shatter. Fragmented strata slide down into the valleys where they continue to break into pieces due to the relentless sun by day and chill by night. Successive strata are exposed to further weathering. The relief of the internal pressure that has built up in rocks that have been underground for aeons can cause them to shatter.[32] Exfoliation also occurs when the outer surfaces of rocks split off in flat flakes. This is believed to be caused by the stresses put on the rock by repeated thermal expansions and contractions which induces fracturing parallel to the original surface.[30] Chemical weathering processes probably play a more important role in deserts than was previously thought. The necessary moisture may be present in the form of dew or mist. Ground water may be drawn to the surface by evaporation and the formation of salt crystals may dislodge rock particles as sand or disintegrate rocks by exfoliation. Shallow caves are sometimes formed at the base of cliffs by this means.[30]
+
+As the desert mountains decay, large areas of shattered rock and rubble occur. The process continues and the end products are either dust or sand. Dust is formed from solidified clay or volcanic deposits whereas sand results from the fragmentation of harder granites, limestone and sandstone.[33] There is a certain critical size (about 0.5 mm) below which further temperature-induced weathering of rocks does not occur and this provides a minimum size for sand grains.[34]
+
+As the mountains are eroded, more and more sand is created. At high wind speeds, sand grains are picked up off the surface and blown along, a process known as saltation. The whirling airborne grains act as a sand blasting mechanism which grinds away solid objects in its path as the kinetic energy of the wind is transferred to the ground.[35] The sand eventually ends up deposited in level areas known as sand-fields or sand-seas, or piled up in dunes.[36]
+
+Sand and dust storms are natural events that occur in arid regions where the land is not protected by a covering of vegetation. Dust storms usually start in desert margins rather than the deserts themselves where the finer materials have already been blown away. As a steady wind begins to blow, fine particles lying on the exposed ground begin to vibrate. At greater wind speeds, some particles are lifted into the air stream. When they land, they strike other particles which may be jerked into the air in their turn, starting a chain reaction. Once ejected, these particles move in one of three possible ways, depending on their size, shape and density; suspension, saltation or creep. Suspension is only possible for particles less than 0.1 mm (0.004 in) in diameter. In a dust storm, these fine particles are lifted up and wafted aloft to heights of up to 6 km (3.7 mi). They reduce visibility and can remain in the atmosphere for days on end, conveyed by the trade winds for distances of up to 6,000 km (3,700 mi).[37] Denser clouds of dust can be formed in stronger winds, moving across the land with a billowing leading edge. The sunlight can be obliterated and it may become as dark as night at ground level.[38] In a study of a dust storm in China in 2001, it was estimated that 6.5 million tons of dust were involved, covering an area of 134,000,000 km2 (52,000,000 sq mi). The mean particle size was 1.44 μm.[39] A much smaller scale, short-lived phenomenon can occur in calm conditions when hot air near the ground rises quickly through a small pocket of cooler, low-pressure air above forming a whirling column of particles, a dust devil.[40]
+
+Sandstorms occur with much less frequency than dust storms. They are often preceded by severe dust storms and occur when the wind velocity increases to a point where it can lift heavier particles. These grains of sand, up to about 0.5 mm (0.020 in) in diameter are jerked into the air but soon fall back to earth, ejecting other particles in the process. Their weight prevents them from being airborne for long and most only travel a distance of a few meters (yards). The sand streams along above the surface of the ground like a fluid, often rising to heights of about 30 cm (12 in).[37] In a really severe steady blow, 2 m (6 ft 7 in) is about as high as the sand stream can rise as the largest sand grains do not become airborne at all. They are transported by creep, being rolled along the desert floor or performing short jumps.[38]
+
+During a sandstorm, the wind-blown sand particles become electrically charged. Such electric fields, which range in size up to 80 kV/m, can produce sparks and cause interference with telecommunications equipment. They are also unpleasant for humans and can cause headaches and nausea.[38] The electric fields are caused by the collision between airborne particles and by the impacts of saltating sand grains landing on the ground. The mechanism is little understood but the particles usually have a negative charge when their diameter is under 250 μm and a positive one when they are over 500 μm.[41][42]
+
+Deserts take up about one third of the Earth's land surface.[8]
+Bottomlands may be salt-covered flats. Eolian processes are major factors in shaping desert landscapes. Polar deserts (also seen as "cold deserts") have similar features, except the main form of precipitation is snow rather than rain. Antarctica is the world's largest cold desert (composed of about 98% thick continental ice sheet and 2% barren rock). Some of the barren rock is to be found in the so-called Dry Valleys of Antarctica that almost never get snow, which can have ice-encrusted saline lakes that suggest evaporation far greater than the rare snowfall due to the strong katabatic winds that even evaporate ice.
+
+Deserts, both hot and cold, play a part in moderating the Earth's temperature. This is because they reflect more of the incoming light and their albedo is higher than that of forests or the sea.[44]
+
+Many people think of deserts as consisting of extensive areas of billowing sand dunes because that is the way they are often depicted on TV and in films,[45] but deserts do not always look like this.[46] Across the world, around 20% of desert is sand, varying from only 2% in North America to 30% in Australia and over 45% in Central Asia.[47] Where sand does occur, it is usually in large quantities in the form of sand sheets or extensive areas of dunes.[47]
+
+A sand sheet is a near-level, firm expanse of partially consolidated particles in a layer that varies from a few centimeters to a few meters thick. The structure of the sheet consists of thin horizontal layers of coarse silt and very fine to medium grain sand, separated by layers of coarse sand and pea-gravel which are a single grain thick. These larger particles anchor the other particles in place and may also be packed together on the surface so as to form a miniature desert pavement.[48] Small ripples form on the sand sheet when the wind exceeds 24 km/h (15 mph). They form perpendicular to the wind direction and gradually move across the surface as the wind continues to blow. The distance between their crests corresponds to the average length of jumps made by particles during saltation. The ripples are ephemeral and a change in wind direction causes them to reorganise.[49]
+
+Sand dunes are accumulations of windblown sand piled up in mounds or ridges. They form downwind of copious sources of dry, loose sand and occur when topographic and climatic conditions cause airborne particles to settle. As the wind blows, saltation and creep take place on the windward side of the dune and individual grains of sand move uphill. When they reach the crest, they cascade down the far side. The upwind slope typically has a gradient of 10° to 20° while the lee slope is around 32°, the angle at which loose dry sand will slip. As this wind-induced movement of sand grains takes place, the dune moves slowly across the surface of the ground.[50] Dunes are sometimes solitary, but they are more often grouped together in dune fields. When these are extensive, they are known as sand seas or ergs.[51]
+
+The shape of the dune depends on the characteristics of the prevailing wind. Barchan dunes are produced by strong winds blowing across a level surface and are crescent-shaped with the concave side away from the wind. When there are two directions from which winds regularly blow, a series of long, linear dunes known as seif dunes may form. These also occur parallel to a strong wind that blows in one general direction. Transverse dunes run at a right angle to the prevailing wind direction. Star dunes are formed by variable winds, and have several ridges and slip faces radiating from a central point. They tend to grow vertically; they can reach a height of 500 m (1,600 ft), making them the tallest type of dune. Rounded mounds of sand without a slip face are the rare dome dunes, found on the upwind edges of sand seas.[51]
+
+A large part of the surface area of the world's deserts consists of flat, stone-covered plains dominated by wind erosion. In "eolian deflation", the wind continually removes fine-grained material, which becomes wind-blown sand. This exposes coarser-grained material, mainly pebbles with some larger stones or cobbles,[36][47] leaving a desert pavement, an area of land overlaid by closely packed smooth stones forming a tessellated mosaic. Different theories exist as to how exactly the pavement is formed. It may be that after the sand and dust is blown away by the wind the stones jiggle themselves into place; alternatively, stones previously below ground may in some way work themselves to the surface. Very little further erosion takes place after the formation of a pavement, and the ground becomes stable. Evaporation brings moisture to the surface by capillary action and calcium salts may be precipitated, binding particles together to form a desert conglomerate.[52] In time, bacteria that live on the surface of the stones accumulate a film of minerals and clay particles, forming a shiny brown coating known as desert varnish.[53]
+
+Other non-sandy deserts consist of exposed outcrops of bedrock, dry soils or aridisols, and a variety of landforms affected by flowing water, such as alluvial fans, sinks or playas, temporary or permanent lakes, and oases.[47] A hamada is a type of desert landscape consisting of a high rocky plateau where the sand has been removed by aeolian processes. Other landforms include plains largely covered by gravels and angular boulders, from which the finer particles have been stripped by the wind. These are called "reg" in the western Sahara, "serir" in the eastern Sahara, "gibber plains" in Australia and "saï" in central Asia.[54] The Tassili Plateau in Algeria is an impressive jumble of eroded sandstone outcrops, canyons, blocks, pinnacles, fissures, slabs and ravines. In some places the wind has carved holes or arches, and in others, it has created mushroom-like pillars narrower at the base than the top.[55] In the Colorado Plateau it is water that has been the eroding force. Here the Colorado River has cut its way over the millennia through the high desert floor creating a canyon that is over a mile (6,000 feet or 1,800 meters) deep in places, exposing strata that are over two billion years old.[56]
+
+One of the driest places on Earth is the Atacama Desert.[57][58][59][60][61] It is virtually devoid of life because it is blocked from receiving precipitation by the Andes mountains to the east and the Chilean Coast Range to the west. The cold Humboldt Current and the anticyclone of the Pacific are essential to keep the dry climate of the Atacama. The average precipitation in the Chilean region of Antofagasta is just 1 mm (0.039 in) per year. Some weather stations in the Atacama have never received rain. Evidence suggests that the Atacama may not have had any significant rainfall from 1570 to 1971. It is so arid that mountains that reach as high as 6,885 m (22,589 ft) are completely free of glaciers and, in the southern part from 25°S to 27°S, may have been glacier-free throughout the Quaternary, though permafrost extends down to an altitude of 4,400 m (14,400 ft) and is continuous above 5,600 m (18,400 ft).[62][63] Nevertheless, there is some plant life in the Atacama, in the form of specialist plants that obtain moisture from dew and the fogs that blow in from the Pacific.[57]
+
+When rain falls in deserts, as it occasionally does, it is often with great violence. The desert surface is evidence of this with dry stream channels known as arroyos or wadis meandering across its surface. These can experience flash floods, becoming raging torrents with surprising rapidity after a storm that may be many kilometers away. Most deserts are in basins with no drainage to the sea but some are crossed by exotic rivers sourced in mountain ranges or other high rainfall areas beyond their borders. The River Nile, the Colorado River and the Yellow River do this, losing much of their water through evaporation as they pass through the desert and raising groundwater levels nearby. There may also be underground sources of water in deserts in the form of springs, aquifers, underground rivers or lakes. Where these lie close to the surface, wells can be dug and oases may form where plant and animal life can flourish.[47] The Nubian Sandstone Aquifer System under the Sahara Desert is the largest known accumulation of fossil water. The Great Man-Made River is a scheme launched by Libya's Muammar Gadaffi to tap this aquifer and supply water to coastal cities.[64] Kharga Oasis in Egypt is 150 km (93 mi) long and is the largest oasis in the Libyan Desert. A lake occupied this depression in ancient times and thick deposits of sandy-clay resulted. Wells are dug to extract water from the porous sandstone that lies underneath.[65] Seepages may occur in the walls of canyons and pools may survive in deep shade near the dried up watercourse below.[66]
+
+Lakes may form in basins where there is sufficient precipitation or meltwater from glaciers above. They are usually shallow and saline, and wind blowing over their surface can cause stress, moving the water over nearby low-lying areas. When the lakes dry up, they leave a crust or hardpan behind. This area of deposited clay, silt or sand is known as a playa. The deserts of North America have more than one hundred playas, many of them relics of Lake Bonneville which covered parts of Utah, Nevada and Idaho during the last ice age when the climate was colder and wetter.[67] These include the Great Salt Lake, Utah Lake, Sevier Lake and many dry lake beds. The smooth flat surfaces of playas have been used for attempted vehicle speed records at Black Rock Desert and Bonneville Speedway and the United States Air Force uses Rogers Dry Lake in the Mojave Desert as runways for aircraft and the space shuttle.[47]
+
+Plants face severe challenges in arid environments. Problems they need to solve include how to obtain enough water, how to avoid being eaten and how to reproduce. Photosynthesis is the key to plant growth. It can only take place during the day as energy from the sun is required, but during the day, many deserts become very hot. Opening stomata to allow in the carbon dioxide necessary for the process causes evapotranspiration, and conservation of water is a top priority for desert vegetation. Some plants have resolved this problem by adopting crassulacean acid metabolism, allowing them to open their stomata during the night to allow CO2 to enter, and close them during the day,[68] or by using C4 carbon fixation.[69]
+
+Many desert plants have reduced the size of their leaves or abandoned them altogether. Cacti are desert specialists, and in most species, the leaves have been dispensed with and the chlorophyll displaced into the trunks, the cellular structure of which has been modified to allow them to store water. When rain falls, the water is rapidly absorbed by the shallow roots and retained to allow them to survive until the next downpour, which may be months or years away.[70] The giant saguaro cacti of the Sonoran Desert form "forests", providing shade for other plants and nesting places for desert birds. Saguaro grows slowly but may live for up to two hundred years. The surface of the trunk is folded like a concertina, allowing it to expand, and a large specimen can hold eight tons of water after a good downpour.[70]
+
+Cacti are present in both North and South America with a post-Gondwana origin. Other xerophytic plants have developed similar strategies by a process known as convergent evolution.[71] They limit water loss by reducing the size and number of stomata, by having waxy coatings and hairy or tiny leaves. Some are deciduous, shedding their leaves in the driest season, and others curl their leaves up to reduce transpiration. Others store water in succulent leaves or stems or in fleshy tubers. Desert plants maximize water uptake by having shallow roots that spread widely, or by developing long taproots that reach down to deep rock strata for ground water.[72] The saltbush in Australia has succulent leaves and secretes salt crystals, enabling it to live in saline areas.[72][73] In common with cacti, many have developed spines to ward off browsing animals.[70]
+
+Some desert plants produce seed which lies dormant in the soil until sparked into growth by rainfall. With annuals, such plants grow with great rapidity and may flower and set seed within weeks, aiming to complete their development before the last vestige of water dries up. For perennial plants, reproduction is more likely to be successful if the seed germinates in a shaded position, but not so close to the parent plant as to be in competition with it. Some seed will not germinate until it has been blown about on the desert floor to scarify the seed coat. The seed of the mesquite tree, which grows in deserts in the Americas, is hard and fails to sprout even when planted carefully. When it has passed through the gut of a pronghorn it germinates readily, and the little pile of moist dung provides an excellent start to life well away from the parent tree.[70]  The stems and leaves of some plants lower the surface velocity of sand-carrying winds and protect the ground from erosion. Even small fungi and microscopic plant organisms found on the soil surface (so-called cryptobiotic soil) can be a vital link in preventing erosion and providing support for other living organisms. Cold deserts often have high concentrations of salt in the soil. Grasses and low shrubs are the dominant vegetation here and the ground may be covered with lichens. Most shrubs have spiny leaves and shed them in the coldest part of the year.[74]
+
+Animals adapted to live in deserts are called xerocoles. There is no evidence that body temperature of mammals and birds is adaptive to the different climates, either of great heat or cold. In fact, with a very few exceptions, their basal metabolic rate is determined by body size, irrespective of the climate in which they live.[75] Many desert animals (and plants) show especially clear evolutionary adaptations for water conservation or heat tolerance and so are often studied in comparative physiology, ecophysiology, and evolutionary physiology. One well-studied example is the specializations of mammalian kidneys shown by desert-inhabiting species.[76] Many examples of convergent evolution have been identified in desert organisms, including between cacti and Euphorbia, kangaroo rats and jerboas, Phrynosoma and Moloch lizards.[77]
+
+Deserts present a very challenging environment for animals. Not only do they require food and water but they also need to keep their body temperature at a tolerable level. In many ways, birds are the ablest to do this of the higher animals. They can move to areas of greater food availability as the desert blooms after local rainfall and can fly to faraway waterholes. In hot deserts, gliding birds can remove themselves from the over-heated desert floor by using thermals to soar in the cooler air at great heights. In order to conserve energy, other desert birds run rather than fly. The cream-colored courser flits gracefully across the ground on its long legs, stopping periodically to snatch up insects. Like other desert birds, it is well-camouflaged by its coloring and can merge into the landscape when stationary. The sandgrouse is an expert at this and nests on the open desert floor dozens of kilometers (miles) away from the waterhole it needs to visit daily. Some small diurnal birds are found in very restricted localities where their plumage matches the color of the underlying surface. The desert lark takes frequent dust baths which ensures that it matches its environment.[78]
+
+Water and carbon dioxide are metabolic end products of oxidation of fats, proteins, and carbohydrates.[79] Oxidising a gram of carbohydrate produces 0.60 grams of water; a gram of protein produces 0.41 grams of water; and a gram of fat produces 1.07 grams of water,[80] making it possible for xerocoles to live with little or no access to drinking water.[81] The kangaroo rat for example makes use of this water of metabolism and conserves water both by having a low basal metabolic rate and by remaining underground during the heat of the day,[82] reducing loss of water through its skin and respiratory system when at rest.[81][83] Herbivorous mammals obtain moisture from the plants they eat. Species such as the addax antelope,[84] dik-dik, Grant's gazelle and oryx are so efficient at doing this that they apparently never need to drink.[85] The camel is a superb example of a mammal adapted to desert life. It minimizes its water loss by producing concentrated urine and dry dung, and is able to lose 40% of its body weight through water loss without dying of dehydration.[86] Carnivores can obtain much of their water needs from the body fluids of their prey.[87] Many other hot desert animals are nocturnal, seeking out shade during the day or dwelling underground in burrows. At depths of more than  50 cm (20 in), these remain at between 30 to 32 °C (86 to 90 °F) regardless of the external temperature.[87] Jerboas, desert rats, kangaroo rats and other small rodents emerge from their burrows at night and so do the foxes, coyotes, jackals and snakes that prey on them. Kangaroos keep cool by increasing their respiration rate, panting, sweating and moistening the skin of their forelegs with saliva.[88] Mammals living in cold deserts have developed greater insulation through warmer body fur and insulating layers of fat beneath the skin. The arctic weasel has a metabolic rate that is two or three times as high as would be expected for an animal of its size. Birds have avoided the problem of losing heat through their feet by not attempting to maintain them at the same temperature as the rest of their bodies, a form of adaptive insulation.[75] The emperor penguin has dense plumage, a downy under layer, an air insulation layer next the skin and various thermoregulatory strategies to maintain its body temperature in one of the harshest environments on Earth.[89]
+
+Being ectotherms, reptiles are unable to live in cold deserts but are well-suited to hot ones. In the heat of the day in the Sahara, the temperature can rise to 50 °C (122 °F). Reptiles cannot survive at this temperature and lizards will be prostrated by heat at 45 °C (113 °F). They have few adaptations to desert life and are unable to cool themselves by sweating so they shelter during the heat of the day. In the first part of the night, as the ground radiates the heat absorbed during the day, they emerge and search for prey. Lizards and snakes are the most numerous in arid regions and certain snakes have developed a novel method of locomotion that enables them to move sidewards and navigate high sand-dunes. These include the horned viper of Africa and the sidewinder of North America, evolutionarily distinct but with similar behavioural patterns because of convergent evolution. Many desert reptiles are ambush predators and often bury themselves in the sand, waiting for prey to come within range.[90]
+
+Amphibians might seem unlikely desert-dwellers, because of their need to keep their skins moist and their dependence on water for reproductive purposes. In fact, the few species that are found in this habitat have made some remarkable adaptations. Most of them are fossorial, spending the hot dry months aestivating in deep burrows. While there they shed their skins a number of times and retain the remnants around them as a waterproof cocoon to retain moisture. In the Sonoran Desert, Couch's spadefoot toad spends most of the year dormant in its burrow. Heavy rain is the trigger for emergence and the first male to find a suitable pool calls to attract others. Eggs are laid and the tadpoles grow rapidly as they must reach metamorphosis before the water evaporates. As the desert dries out, the adult toads rebury themselves. The juveniles stay on the surface for a while, feeding and growing, but soon dig themselves burrows. Few make it to adulthood.[91] The water holding frog in Australia has a similar life cycle and may aestivate for as long as five years if no rain falls.[92] The Desert rain frog of Namibia is nocturnal and survives because of the damp sea fogs that roll in from the Atlantic.[93]
+
+Invertebrates, particularly arthropods, have successfully made their homes in the desert. Flies, beetles, ants, termites, locusts, millipedes, scorpions and spiders[94] have hard cuticles which are impervious to water and many of them lay their eggs underground and their young develop away from the temperature extremes at the surface.[95] The Saharan silver ant (Cataglyphis bombycina) uses a heat shock protein in a novel way and forages in the open during brief forays in the heat of the day.[96] The long-legged darkling beetle in Namibia stands on its front legs and raises its carapace to catch the morning mist as condensate, funnelling the water into its mouth.[97] Some arthropods make use of the ephemeral pools that form after rain and complete their life cycle in a matter of days. The desert shrimp does this, appearing "miraculously" in new-formed puddles as the dormant eggs hatch. Others, such as brine shrimps, fairy shrimps and tadpole shrimps, are cryptobiotic and can lose up to 92% of their bodyweight, rehydrating as soon as it rains and their temporary pools reappear.[98]
+
+Humans have long made use of deserts as places to live,[99] and more recently have started to exploit them for minerals[100] and energy capture.[101] Deserts play a significant role in human culture with an extensive literature.[102]
+
+People have been living in deserts for millennia. Many, such as the Bushmen in the Kalahari, the Aborigines in Australia and various tribes of North American Indians, were originally hunter-gatherers. They developed skills in the manufacture and use of weapons, animal tracking, finding water, foraging for edible plants and using the things they found in their natural environment to supply their everyday needs. Their self-sufficient skills and knowledge were passed down through the generations by word of mouth.[99] Other cultures developed a nomadic way of life as herders of sheep, goats, cattle, camels, yaks, llamas or reindeer. They travelled over large areas with their herds, moving to new pastures as seasonal and erratic rainfall encouraged new plant growth. They took with them their tents made of cloth or skins draped over poles and their diet included milk, blood and sometimes meat.[103]
+
+The desert nomads were also traders. The Sahara is a very large expanse of land stretching from the Atlantic rim to Egypt. Trade routes were developed linking the Sahel in the south with the fertile Mediterranean region to the north and large numbers of camels were used to carry valuable goods across the desert interior. The Tuareg were traders and the goods transported traditionally included slaves, ivory and gold going northwards and salt going southwards. Berbers with knowledge of the region were employed to guide the caravans between the various oases and wells.[104] Several million slaves may have been taken northwards across the Sahara between the 8th and 18th centuries.[105] Traditional means of overland transport declined with the advent of motor vehicles, shipping and air freight, but caravans still travel along routes between Agadez and Bilma and between Timbuktu and Taoudenni carrying salt from the interior to desert-edge communities.[106]
+
+Round the rims of deserts, where more precipitation occurred and conditions were more suitable, some groups took to cultivating crops. This may have happened when drought caused the death of herd animals, forcing herdsmen to turn to cultivation. With few inputs, they were at the mercy of the weather and may have lived at bare subsistence level. The land they cultivated reduced the area available to nomadic herders, causing disputes over land. The semi-arid fringes of the desert have fragile soils which are at risk of erosion when exposed, as happened in the American Dust Bowl in the 1930s. The grasses that held the soil in place were ploughed under, and a series of dry years caused crop failures, while enormous dust storms blew the topsoil away. Half a million Americans were forced to leave their land in this catastrophe.[107]
+
+Similar damage is being done today to the semi-arid areas that rim deserts and about twelve million hectares of land are being turned to desert each year.[108] Desertification is caused by such factors as drought, climatic shifts, tillage for agriculture, overgrazing and deforestation. Vegetation plays a major role in determining the composition of the soil. In many environments, the rate of erosion and run off increases dramatically with reduced vegetation cover.[109]
+
+Deserts contain substantial mineral resources, sometimes over their entire surface, giving them their characteristic colors. For example, the red of many sand deserts comes from laterite minerals.[110] Geological processes in a desert climate can concentrate minerals into valuable deposits. Leaching by ground water can extract ore minerals and redeposit them, according to the water table, in concentrated form.[100] Similarly, evaporation tends to concentrate minerals in desert lakes, creating dry lake beds or playas rich in minerals. Evaporation can concentrate minerals as a variety of evaporite deposits, including gypsum, sodium nitrate, sodium chloride and borates.[100] Evaporites are found in the USA's Great Basin Desert, historically exploited by the "20-mule teams" pulling carts of borax from Death Valley to the nearest railway.[100] A desert especially rich in mineral salts is the Atacama Desert, Chile, where sodium nitrate has been mined for explosives and fertilizer since around 1850.[100] Other desert minerals are copper from Chile, Peru, and Iran, and iron and uranium in Australia. Many other metals, salts and commercially valuable types of rock such as pumice are extracted from deserts around the world.[100]
+
+Oil and gas form on the bottom of shallow seas when micro-organisms decompose under anoxic conditions and later become covered with sediment. Many deserts were at one time the sites of shallow seas and others have had underlying hydrocarbon deposits transported to them by the movement of tectonic plates.[111]
+Some major oilfields such as Ghawar are found under the sands of Saudi Arabia.[100] Geologists believe that other oil deposits were formed by aeolian processes in ancient deserts as may be the case with some of the major American oil fields.[100]
+
+Traditional desert farming systems have long been established in North Africa, irrigation being the key to success in an area where water stress is a limiting factor to growth. Techniques that can be used include drip irrigation, the use of organic residues or animal manures as fertilisers and other traditional agricultural management practices. Once fertility has been built up, further crop production preserves the soil from destruction by wind and other forms of erosion.[112] It has been found that plant growth-promoting bacteria play a role in increasing the resistance of plants to stress conditions and these rhizobacterial suspensions could be inoculated into the soil in the vicinity of the plants. A study of these microbes found that desert farming hampers desertification by establishing islands of fertility allowing farmers to achieve increased yields despite the adverse environmental conditions.[112] A field trial in the Sonoran Desert which exposed the roots of different species of tree to rhizobacteria and the nitrogen fixing bacterium Azospirillum brasilense with the aim of restoring degraded lands was only partially successful.[112]
+
+The Judean Desert was farmed in the 7th century BC during the Iron Age to supply food for desert forts.[113] Native Americans in the south western United States became agriculturalists around 600 AD when seeds and technologies became available from Mexico. They used terracing techniques and grew gardens beside seeps, in moist areas at the foot of dunes, near streams providing flood irrigation and in areas irrigated by extensive specially built canals. The Hohokam tribe constructed over 500 miles (800 km) of large canals and maintained them for centuries, an impressive feat of engineering. They grew maize, beans, squash and peppers.[114]
+
+A modern example of desert farming is the Imperial Valley in California, which has high temperatures and average rainfall of just 3 in (76 mm) per year.[115]  The economy is heavily based on agriculture and the land is irrigated through a network of canals and pipelines sourced entirely from the Colorado River via the All-American Canal. The soil is deep and fertile, being part of the river's flood plains, and what would otherwise have been desert has been transformed into one of the most productive farming regions in California. Other water from the river is piped to urban communities but all this has been at the expense of the river, which below the extraction sites no longer has any above-ground flow during most of the year. Another problem of growing crops in this way is the build-up of salinity in the soil caused by the evaporation of river water.[116] The greening of the desert remains an aspiration and was at one time viewed as a future means for increasing food production for the world's growing population. This prospect has proved false as it disregarded the environmental damage caused elsewhere by the diversion of water for desert project irrigation.[117]
+
+Deserts are increasingly seen as sources for solar energy, partly due to low amounts of cloud cover. Many solar power plants have been built in the Mojave Desert such as the Solar Energy Generating Systems and Ivanpah Solar Power Facility.[118] Large swaths of this desert are covered in mirrors.[119]
+
+The potential for generating solar energy from the Sahara Desert is huge, the highest found on the globe. Professor David Faiman of Ben-Gurion University has stated that the technology now exists to supply all of the world's electricity needs from 10% of the Sahara Desert.[120] Desertec Industrial Initiative was a consortium seeking $560 billion to invest in North African solar and wind installations over the next forty years to supply electricity to Europe via cable lines running under the Mediterranean Sea. European interest in the Sahara Desert stems from its two aspects: the almost continual daytime sunshine and plenty of unused land. The Sahara receives more sunshine per acre than any part of Europe. The Sahara Desert also has the empty space totalling hundreds of square miles required to house fields of mirrors for solar plants.[121]
+
+The Negev Desert, Israel, and the surrounding area, including the Arava Valley, receive plenty of sunshine and are generally not arable. This has resulted in the construction of many solar plants.[101] David Faiman has proposed that "giant" solar plants in the Negev could supply all of Israel's needs for electricity.[120]
+
+The Arabs were probably the first organized force to conduct successful battles in the desert. By knowing back routes and the locations of oases and by utilizing camels, Muslim Arab forces were able to successfully overcome both Roman and Persian forces in the period 600 to 700 AD during the expansion of the Islamic caliphate.[122]
+
+Many centuries later, both world wars saw fighting in the desert. In the First World War, the Ottoman Turks were engaged with the British regular army in a campaign that spanned the Arabian peninsula. The Turks were defeated by the British, who had the backing of irregular Arab forces that were seeking to revolt against the Turks in the Hejaz, made famous in T.E. Lawrence's book Seven Pillars of Wisdom.[123][124]
+
+In the Second World War, the Western Desert Campaign began in Italian Libya. Warfare in the desert offered great scope for tacticians to use the large open spaces without the distractions of casualties among civilian populations. Tanks and armoured vehicles were able to travel large distances unimpeded and land mines were laid in large numbers. However, the size and harshness of the terrain meant that all supplies needed to be brought in from great distances. The victors in a battle would advance and their supply chain would necessarily become longer, while the defeated army could retreat, regroup and resupply. For these reasons, the front line moved back and forth through hundreds of kilometers as each side lost and regained momentum.[125]  Its most easterly point was at El Alamein in Egypt, where the Allies decisively defeated the Axis forces in 1942.[126]
+
+The desert is generally thought of as a barren and empty landscape. It has been portrayed by writers, film-makers, philosophers, artists and critics as a place of extremes, a metaphor for anything from death, war or religion to the primitive past or the desolate future.[127]
+
+There is an extensive literature on the subject of deserts.[102] An early historical account is that of Marco Polo (c. 1254–1324), who travelled through Central Asia to China, crossing a number of deserts in his twenty four year trek.[128] Some accounts give vivid descriptions of desert conditions, though often accounts of journeys across deserts are interwoven with reflection, as is the case in Charles Montagu Doughty's major work, Travels in Arabia Deserta (1888).[129] Antoine de Saint-Exupéry described both his flying and the desert in Wind, Sand and Stars[130] and Gertrude Bell travelled extensively in the Arabian desert in the early part of the 20th century, becoming an expert on the subject, writing books and advising the British government on dealing with the Arabs.[131] Another woman explorer was Freya Stark who travelled alone in the Middle East, visiting Turkey, Arabia, Yemen, Syria, Persia and Afghanistan, writing over twenty books on her experiences.[132] The German naturalist Uwe George spent several years living in deserts, recording his experiences and research in his book, In the Deserts of this Earth.[133]
+
+The American poet Robert Frost expressed his bleak thoughts in his poem, Desert Places, which ends with the stanza "They cannot scare me with their empty spaces / Between stars – on stars where no human race is. / I have it in me so much nearer home / To scare myself with my own desert places."[134]
+
+Mars is the only other planet in the Solar System besides Earth on which deserts have been identified.[135] Despite its low surface atmospheric pressure (only 1/100 of that of the Earth), the patterns of atmospheric circulation on Mars have formed a sea of circumpolar sand more than 5 million km2 (1.9 million sq mi) in the area, larger than most deserts on Earth. The Martian deserts principally consist of dunes in the form of half-moons in flat areas near the permanent polar ice caps in the north of the planet. The smaller dune fields occupy the bottom of many of the craters situated in the Martian polar regions.[136] Examination of the surface of rocks by laser beamed from the Mars Exploration Rover have shown a surface film that resembles the desert varnish found on Earth although it might just be surface dust.[137] The surface of Titan, a moon of Saturn, also has a desert-like surface with dune seas.[138]

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+The Mesozoic Era ( /ˌmɛz.əˈzoʊ.ɪk, ˌmɛz.oʊ-, ˌmɛs-, ˌmiː.zə-, -zoʊ-, ˌmiː.sə-, -soʊ-/ mez-ə-ZOH-ik, mez-oh-, mess-, mee-zə-, -⁠zoh-, mee-sə-, -⁠soh-)[1][2] is an interval of geological time from about 252 to 66 million years ago. It is also called the Age of Reptiles and the Age of Conifers.[3]
+
+The Mesozoic ("middle life") is one of three geologic eras of the Phanerozoic Eon, preceded by the Paleozoic ("ancient life") and succeeded by the Cenozoic ("new life"). The era is subdivided into three major periods: the Triassic, Jurassic, and Cretaceous, which are further subdivided into a number of epochs and stages.
+
+The era began in the wake of the Permian–Triassic extinction event, the largest well-documented mass extinction in Earth's history, and ended with the Cretaceous–Paleogene extinction event, another mass extinction whose victims included the non-avian dinosaurs. The Mesozoic was a time of significant tectonic, climate, and evolutionary activity. The era witnessed the gradual rifting of the supercontinent Pangaea into separate landmasses that would move into their current positions during the next era. The climate of the Mesozoic was varied, alternating between warming and cooling periods. Overall, however, the Earth was hotter than it is today. Dinosaurs first appeared in the Mid-Triassic, and became the dominant terrestrial vertebrates in the Late Triassic or Early Jurassic, occupying this position for about 150 or 135 million years until their demise at the end of the Cretaceous. Birds first appeared in the Jurassic (however, true toothless birds appeared first in the Cretaceous), having evolved from a branch of theropod dinosaurs. The first mammals also appeared during the Mesozoic, but would remain small—less than 15 kg (33 lb)—until the Cenozoic. The flowering plants (angiosperms) arose in the Triassic or Jurassic and came to prominence in the late Cretaceous when they replaced the conifers and other gymnosperms as the dominant trees.
+
+The phrase "Age of Reptiles" was introduced by the 19th century paleontologist Gideon Mantell who viewed it as dominated by diapsids such as Iguanodon, Megalosaurus, Plesiosaurus, and Pterodactylus.
+
+Mesozoic means "middle life", deriving from the Greek prefix meso-/μεσο- for "between" and zōon/ζῷον meaning "animal" or "living being". The name "Mesozoic" was proposed in 1840 by the British geologist John Phillips (1800–1874).[4][5]
+
+The Mesozoic era was originally described as the "secondary" era, following the primary or Paleozoic, and preceding the Tertiary.[6]
+
+Following the Paleozoic, the Mesozoic extended roughly 186 million years, from 251.902 to 66 million years ago when the Cenozoic Era began. This time frame is separated into three geologic periods. From oldest to youngest:
+
+The lower boundary of the Mesozoic is set by the Permian–Triassic extinction event, during which approximately 90% to 96% of marine species and 70% of terrestrial vertebrates became extinct.[7] It is also known as the "Great Dying" because it is considered the largest mass extinction in the Earth's history. The upper boundary of the Mesozoic is set at the Cretaceous–Paleogene extinction event (or K–Pg extinction event[8]), which may have been caused by an asteroid impactor that created Chicxulub Crater on the Yucatán Peninsula. Towards the Late Cretaceous, large volcanic eruptions are also believed to have contributed to the Cretaceous–Paleogene extinction event. Approximately 50% of all genera became extinct, including all of the non-avian dinosaurs.
+
+The Triassic ranges roughly from 252 million to 201 million years ago, preceding the Jurassic Period. The period is bracketed between the Permian–Triassic extinction event and the Triassic–Jurassic extinction event, two of the "big five", and it is divided into three major epochs: Early, Middle, and Late Triassic.[9]
+
+The Early Triassic, about 252 to 247 million years ago, was dominated by deserts in the interior of the Pangaea supercontinent. The Earth had just witnessed a massive die-off in which 95% of all life became extinct, and the most common vertebrate life on land were Lystrosaurus, labyrinthodonts, and Euparkeria along with many other creatures that managed to survive the Permian extinction. Temnospondyls evolved during this time and would be the dominant predator for much of the Triassic.[10]
+
+The Middle Triassic, from 247 to 237 million years ago, featured the beginnings of the breakup of Pangaea and the opening of the Tethys Ocean. Ecosystems had recovered from the Permian extinction. Algae, sponge, corals, and crustaceans all had recovered, and new aquatic reptiles evolved, such as ichthyosaurs and nothosaurs. On land, pine forests flourished, as did groups of insects like mosquitoes and fruit flies. Reptiles began to get bigger and bigger, and the first crocodilians and dinosaurs evolved, which sparked competition with the large amphibians that had previously ruled the freshwater world, respectively mammal-like reptiles on land.[11]
+
+Following the bloom of the Middle Triassic, the Late Triassic, from 237 to 201 million years ago, featured frequent heat spells and moderate precipitation (10–20 inches per year). The recent warming led to a boom of dinosaurian evolution on land as those one began to separate from each other (Nyasasaurus from 243 to 210 million years ago, approximately 235–30 ma, some of them separated into Sauropodomorphs, Theropods and Herrerasaurids), as well as the first pterosaurs. During the Late Triassic, some advanced cynodonts gave rise to the first Mammaliaformes. All this climatic change, however, resulted in a large die-out known as the Triassic–Jurassic extinction event, in which many archosaurs (excluding pterosaurs, dinosaurs and crocodylomorphs), most synapsids, and almost all large amphibians became extinct, as well as 34% of marine life, in the Earth's fourth mass extinction event. The cause is debatable;[12][13] flood basalt eruptions at the Central Atlantic magmatic province is cited as one possible cause.
+
+The Jurassic ranges from 200 million years to 145 million years ago and features three major epochs: The Early Jurassic, the Middle Jurassic, and the Late Jurassic.[14]
+
+The Early Jurassic spans from 200 to 175 million years ago.[14] The climate was tropical, much more humid than the Triassic. In the oceans, plesiosaurs, ichthyosaurs and ammonites were abundant. On land, dinosaurs and other archosaurs staked their claim as the dominant race, with theropods such as Dilophosaurus at the top of the food chain. The first true crocodiles evolved, pushing the large amphibians to near extinction. All-in-all, archosaurs rose to rule the world. Meanwhile, the first true mammals evolved, remaining relatively small but spreading widely; the Jurassic Castorocauda, for example, had adaptations for swimming, digging and catching fish. Fruitafossor, from the late Jurassic period about 150 million years ago, was about the size of a chipmunk, and its teeth, forelimbs and back suggest that it dug open the nests of social insects (probably termites, as ants had not yet appeared). The first multituberculates like Rugosodon evolved, while volaticotherians took to the skies.
+
+The Middle Jurassic spans from 175 to 163 million years ago.[14] During this epoch, dinosaurs flourished as huge herds of sauropods, such as Brachiosaurus and Diplodocus, filled the fern prairies, chased by many new predators such as Allosaurus. Conifer forests made up a large portion of the forests. In the oceans, plesiosaurs were quite common, and ichthyosaurs flourished. This epoch was the peak of the reptiles.[15]
+
+The Late Jurassic spans from 163 to 145 million years ago.[14] During this epoch, the first avialans, like Archaeopteryx, evolved from small coelurosaurian dinosaurs. The increase in sea levels opened up the Atlantic seaway, which has grown continually larger until today. The divided landmasses gave opportunity for the diversification of new dinosaurs.
+
+The Cretaceous is the longest period of the Mesozoic, but has only two epochs: Early and Late Cretaceous.[16]
+
+The Early Cretaceous spans from 145 to 100 million years ago.[16] The Early Cretaceous saw the expansion of seaways, and as a result, the decline and/or extinction of Laurasian sauropods. Some island-hopping dinosaurs, like Eustreptospondylus, evolved to cope with the coastal shallows and small islands of ancient Europe. Other dinosaurs rose up to fill the empty space that the Jurassic-Cretaceous extinction left behind, such as Carcharodontosaurus and Spinosaurus. Of the most successful was the Iguanodon, which spread to every continent. Seasons came back into effect and the poles got seasonally colder, but some dinosaurs still inhabited the polar forests year round, such as Leaellynasaura and Muttaburrasaurus. The poles were too cold for crocodiles, and became the last stronghold for large amphibians like Koolasuchus. Pterosaurs got larger as genera like Tapejara and Ornithocheirus evolved. Mammals continued to expand their range: eutriconodonts produced fairly large, wolverine-like predators like Repenomamus and Gobiconodon, early therians began to expand into metatherians and eutherians, and cimolodont multituberculates went on to become common in the fossil record.
+
+The Late Cretaceous spans from 100 to 66 million years ago. The Late Cretaceous featured a cooling trend that would continue in the Cenozoic era. Eventually, tropics were restricted to the equator and areas beyond the tropic lines experienced extreme seasonal changes in weather. Dinosaurs still thrived, as new taxa such as Tyrannosaurus, Ankylosaurus, Triceratops and hadrosaurs dominated the food web. In the oceans, mosasaurs ruled, filling the role of the ichthyosaurs, which, after declining, had disappeared in the Cenomanian-Turonian boundary event. Though pliosaurs had gone extinct in the same event, long-necked plesiosaurs such as Elasmosaurus continued to thrive. Flowering plants, possibly appearing as far back as the Triassic, became truly dominant for the first time. Pterosaurs in the Late Cretaceous declined for poorly understood reasons, though this might be due to tendencies of the fossil record, as their diversity seems to be much higher than previously thought. Birds became increasingly common and diversified into a variety of enantiornithe and ornithurine forms. Though mostly small, marine hesperornithes became relatively large and flightless, adapted to life in the open sea. Metatherians and primitive eutherian also became common and even produced large and specialised genera like Didelphodon and Schowalteria. Still, the dominant mammals were multituberculates, cimolodonts in the north and gondwanatheres in the south. At the end of the Cretaceous, the Deccan traps and other volcanic eruptions were poisoning the atmosphere. As this continued, it is thought that a large meteor smashed into earth 66 million years ago, creating the Chicxulub Crater in an event known as the K-Pg Extinction (formerly K-T), the fifth and most recent mass extinction event, in which 75% of life became extinct, including all non-avian dinosaurs.[17] Everything over 10 kilograms became extinct. The age of the dinosaurs was over.[18][19]
+
+Compared to the vigorous convergent plate mountain-building of the late Paleozoic, Mesozoic tectonic deformation was comparatively mild. The sole major Mesozoic orogeny occurred in what is now the Arctic, creating the Innuitian orogeny, the Brooks Range, the Verkhoyansk and Cherskiy Ranges in Siberia, and the Khingan Mountains in Manchuria.
+
+This orogeny was related to the opening of the Arctic Ocean and subduction of the North China and Siberian cratons under the Pacific Ocean.[20] In contrast, the era featured the dramatic rifting of the supercontinent Pangaea, which gradually split into a northern continent, Laurasia, and a southern continent, Gondwana. This created the passive continental margin that characterizes most of the Atlantic coastline (such as along the U.S. East Coast) today.[21]
+
+By the end of the era, the continents had rifted into nearly their present forms, though not their present positions. Laurasia became North America and Eurasia, while Gondwana split into South America, Africa, Australia, Antarctica and the Indian subcontinent, which collided with the Asian plate during the Cenozoic, giving rise to the Himalayas.
+
+The Triassic was generally dry, a trend that began in the late Carboniferous, and highly seasonal, especially in the interior of Pangaea. Low sea levels may have also exacerbated temperature extremes. With its high specific heat capacity, water acts as a temperature-stabilizing heat reservoir, and land areas near large bodies of water—especially oceans—experience less variation in temperature. Because much of Pangaea's land was distant from its shores, temperatures fluctuated greatly, and the interior probably included expansive deserts. Abundant red beds and evaporites such as halite support these conclusions, but some evidence suggests the generally dry climate of was punctuated by episodes of increased rainfall.[22] The most important humid episodes were the Carnian Pluvial Event and one in the Rhaetian, a few million years before the Triassic–Jurassic extinction event.
+
+Sea levels began to rise during the Jurassic, probably caused by an increase in seafloor spreading. The formation of new crust beneath the surface displaced ocean waters by as much as 200 m (656 ft) above today's sea level, flooding coastal areas. Furthermore, Pangaea began to rift into smaller divisions, creating new shoreline around the Tethys Ocean. Temperatures continued to increase, then began to stabilize. Humidity also increased with the proximity of water, and deserts retreated.
+
+The climate of the Cretaceous is less certain and more widely disputed. Probably, higher levels of carbon dioxide in the atmosphere are thought to have almost eliminated the north–south temperature gradient: temperatures were about the same across the planet, and about 10°C higher than today.  The circulation of oxygen to the deep ocean may also have been disrupted,[16][dubious  – discuss] preventing the decomposition of large volumes of organic matter, which was eventually deposited as "black shale".
+
+Not all data support these hypotheses, however. Even with the overall warmth, temperature fluctuations should have been sufficient for the presence of polar ice caps and glaciers, but there is no evidence of either. Quantitative models have also been unable to recreate the flatness of the Cretaceous temperature gradient.[citation needed]
+
+Different studies have come to different conclusions about the amount of oxygen in the atmosphere during different parts of the Mesozoic, with some concluding oxygen levels were lower than the current level (about 21%) throughout the Mesozoic,[23][24] some concluding they were lower in the Triassic and part of the Jurassic but higher in the Cretaceous,[25][26][27] and some concluding they were higher throughout most or all of the Triassic, Jurassic and Cretaceous.[28][29]
+
+The dominant land plant species of the time were gymnosperms, which are vascular, cone-bearing, non-flowering plants such as conifers that produce seeds without a coating. This is opposed to the earth's current flora, in which the dominant land plants in terms of number of species are angiosperms. One particular plant genus, Ginkgo, is thought to have evolved at this time and is represented today by a single species, Ginkgo biloba. As well, the extant genus Sequoia is believed to have evolved in the Mesozoic.[30]
+
+Flowering plants radiated sometime in the early Cretaceous, first in the tropics, but the even temperature gradient allowed them to spread toward the poles throughout the period. By the end of the Cretaceous, angiosperms dominated tree floras in many areas, although some evidence suggests that biomass was still dominated by cycads and ferns until after the Cretaceous–Paleogene extinction. Some plant species had distributions that were markedly different from succeeding periods; for example, the Schizeales, a fern order, were skewed to the Northern Hemisphere in the Mesozoic, but are now better represented in the Southern Hemisphere.[31]
+
+The extinction of nearly all animal species at the end of the Permian Period allowed for the radiation of many new lifeforms.  In particular, the extinction of the large herbivorous pareiasaurs and carnivorous gorgonopsians left those ecological niches empty. Some were filled by the surviving cynodonts and dicynodonts, the latter of which subsequently became extinct.
+
+Recent research indicates that it took much longer for the reestablishment of complex ecosystems with high biodiversity, complex food webs, and specialized animals in a variety of niches, beginning in the mid-Triassic 4M to 6M years after the extinction,[32] and not fully proliferated until 30M years after the extinction.[33] Animal life was then dominated by various archosaurs: dinosaurs, pterosaurs, and aquatic reptiles such as ichthyosaurs, plesiosaurs, and mosasaurs.
+
+The climatic changes of the late Jurassic and Cretaceous favored further adaptive radiation. The Jurassic was the height of archosaur diversity, and the first birds and eutherian mammals also appeared. Some have argued that insects diversified in symbiosis with angiosperms, because insect anatomy, especially the mouth parts, seems particularly well-suited for flowering plants. However, all major insect mouth parts preceded angiosperms, and insect diversification actually slowed when they arrived, so their anatomy originally must have been suited for some other purpose.

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+Yerevan (UK: /ˌjɛrəˈvæn/ YERR-ə-VAN, US: /-ˈvɑːn/ -⁠VAHN; Armenian: Երևան[a] [jɛɾɛˈvɑn] (listen), sometimes spelled Erevan)[b] is the capital and largest city of Armenia and one of the world's oldest continuously inhabited cities.[18] Situated along the Hrazdan River, Yerevan is the administrative, cultural, and industrial center of the country. It has been the capital since 1918, the fourteenth in the history of Armenia and the seventh located in or around the Ararat plain. The city also serves as the seat of the Araratian Pontifical Diocese; the largest diocese of the Armenian Apostolic Church and one of the oldest dioceses in the world.[19]
+
+The history of Yerevan dates back to the 8th century BC, with the founding of the fortress of Erebuni in 782 BC by king Argishti I at the western extreme of the Ararat plain.[20] Erebuni was "designed as a great administrative and religious centre, a fully royal capital."[21] By the late ancient Armenian Kingdom, new capital cities were established and Yerevan declined in importance. Under Iranian and Russian rule, it was the center of the Erivan Khanate from 1736 to 1828 and the Erivan Governorate from 1850 to 1917, respectively. After World War I, Yerevan became the capital of the First Republic of Armenia as thousands of survivors of the Armenian Genocide in the Ottoman Empire arrived in the area.[22] The city expanded rapidly during the 20th century as Armenia became part of the Soviet Union. In a few decades, Yerevan was transformed from a provincial town within the Russian Empire to Armenia's principal cultural, artistic, and industrial center, as well as becoming the seat of national government.
+
+With the growth of the Armenian economy, Yerevan has undergone major transformation. Much construction has been done throughout the city since the early 2000s, and retail outlets such as restaurants, shops, and street cafés, which were rare during Soviet times, have multiplied. As of 2011[update], the population of Yerevan was 1,060,138, just over 35% of the Republic of Armenia's total population. According to the official estimate of 2016, the current population of the city is 1,073,700.[23] Yerevan was named the 2012 World Book Capital by UNESCO.[24] Yerevan is an associate member of Eurocities.[25]
+
+Of the notable landmarks of Yerevan, Erebuni Fortress is considered to be the birthplace of the city, the Katoghike Tsiranavor church is the oldest surviving church of Yerevan and Saint Gregory Cathedral is the largest Armenian cathedral in the world, Tsitsernakaberd is the official memorial to the victims of the Armenian Genocide, and several opera houses, theatres, museums, libraries, and other cultural institutions. Yerevan Opera Theatre is the main spectacle hall of the Armenian capital, the National Gallery of Armenia is the largest art museum in Armenia and shares a building with the History Museum of Armenia, and the Matenadaran repository contains one of the largest depositories of ancient books and manuscripts in the world.
+
+One theory regarding the origin of Yerevan's name is the city was named after the Armenian king, Yervand (Orontes) IV, the last leader of the Orontid Dynasty, and founder of the city of Yervandashat.[27] However, it is likely that the city's name is derived from the Urartian military fortress of Erebuni (Էրեբունի), which was founded on the territory of modern-day Yerevan in 782 BC by Argishti I.[27] As elements of the Urartian language blended with that of the Armenian one, the name eventually evolved into Yerevan (Erebuni = Erevani = Erevan = Yerevan). Scholar Margarit Israelyan notes these changes when comparing inscriptions found on two cuneiform tablets at Erebuni:
+
+The transcription of the second cuneiform bu [original emphasis] of the word was very essential in our interpretation as it is the Urartaean b that has been shifted to the Armenian v (b > v). The original writing of the inscription read «er-bu-ni»; therefore the prominent Armenianologist-orientalist Prof. G. A. Ghapantsian justly objected, remarking that the Urartu b changed to v at the beginning of the word (Biani > Van) or between two vowels (ebani > avan, Zabaha > Javakhk)....In other words b was placed between two vowels. The true pronunciation of the fortress-city was apparently Erebuny.[28]
+
+Early Christian Armenian chroniclers attributed the origin of the name Yerevan to a derivation from an expression exclaimed by Noah, in Armenian. While looking in the direction of Yerevan, after the ark had landed on Mount Ararat and the flood waters had receded, Noah is believed to have exclaimed, "Yerevats!" ("it appeared!").[27]
+
+In the late medieval and early modern periods, when Yerevan was under Turkic and later Persian rule, the city was known in Persian as Iravân (Persian: ایروان‎). This name is still widely used by Azerbaijanis (Azerbaijani: İrəvan). The city was officially known as Erivan (Russian: Эривань) under Russian rule during the 19th and early 20th centuries. The city was renamed back to Yerevan (Ереван) in 1936.[29] Up until the mid-1970s the city's name was spelled Erevan, more often than Yerevan, in English sources.[30][31]
+
+The principal symbol of Yerevan is Mount Ararat, which is visible from any area in the capital. The seal of the city is a crowned lion on a pedestal with the inscriptit in the upper part. The emblem is a rectangular shield with a blue border.[34]
+
+On 27 September 2004, Yerevan adopted an anthem, "Erebuni-Yerevan", written by Paruyr Sevak and composed by Edgar Hovhanisyan. It was selected in a competition for a new anthem and new flag that would best represent the city. The chosen flag has a white background with the city's seal in the middle, surrounded by twelve small red triangles that symbolize the twelve historic capitals of Armenia. The flag includes the three colours of the Armenian National flag. The lion is portrayed on the orange background with blue edging.[35]
+
+The territory of Yerevan has been inhabited since approximately the 2nd half of the 4th millennium BC. The southern part of the city currently known as Shengavit has been populated since at least 3200 BC, during the period of Kura–Araxes culture of the early Bronze Age. The first excavations at the Shengavit historical site was conducted between 1936 and 1938 under the guidance of archaeologist Yevgeny Bayburdyan. After two decades, archaeologist Sandro Sardarian resumed the excavations starting from 1958 until 1983.[36] The 3rd phase of the excavations started in 2000, under the guidance of archaeologist Hakob Simonyan. In 2009, Simonyan was joined by professor Mitchell S. Rothman from the Widener University of Pennsylvania. Together they conducted three series of excavations in 2009, 2010, and 2012 respectively. During the process, a full stratigraphic column to bedrock was reached, showing there to be 8 or 9 distinct stratigraphic levels. These levels cover a time between 3200 BC and 2500 BC. Evidences of later use of the site, possibly until 2200 BC, were also found. The excavation process revealed a series of large round buildings with square adjoining rooms and minor round buildings. A series of ritual installations was discovered in 2010 and 2012.
+
+The ancient kingdom of Urartu was formed in the 9th century BC by King Arame in the basin of Lake Van of the Armenian Highland, including the territory of modern-day Yerevan.[37] Archaeological evidence, such as a cuneiform inscription,[38] indicates that the Urartian military fortress of Erebuni (Էրեբունի) was founded in 782 BC by the orders of King Argishti I at the site of modern-day Yerevan, to serve as a fort and citadel guarding against attacks from the north Caucasus.[27] The cuneiform inscription found at Erebuni Fortress reads:
+
+By the greatness of the God Khaldi, Argishti, son of Menua, built this mighty stronghold and proclaimed it Erebuni for the glory of Biainili [Urartu] and to instill fear among the king's enemies. Argishti says, "The land was a desert, before the great works I accomplished upon it. By the greatness of Khaldi, Argishti, son of Menua, is a mighty king, king of Biainili, and ruler of Tushpa." [Van].[39]
+
+During the height of the Urartian power, irrigation canals and artificial reservoirs were built in Erebuni and its surrounding territories.
+
+In the mid-7th century BC, the city of Teishebaini was built by Rusa II of Urartu, around 7 kilometres (4.3 miles) west of Erebuni Fortress.[40] It was fortified on a hill -currently known as Karmir Blur within Shengavit District of Yerevan- to protect the eastern borders of Urartu from the barbaric Cimmerians and Scythians. During excavations, the remains of a governors palace that contained a hundred and twenty rooms spreading across more than 40,000 m2 (10 acres) was found, along with a citadel dedicated to the Urartian god Teisheba. The construction of the city of Teishebaini, as well as the palace and the citadel was completed by the end of the 7th century BC, during the reign of Rusa III. However, Teishebaini was destroyed by an alliance of Medes and the Scythians in 585 BC.
+
+In 590 BC, following the fall of the Kingdom of Urartu by the hands of the Iranian Medes, Erebuni along with the Armenian Highland became part of the Median Empire.
+
+However, in 550 BC, the Median Empire was conquered by Cyrus the Great, and Erebuni became part of the Achaemenid Empire. Between 522 BC and 331 BC, Erebuni was one of the main centers of the Satrapy of Armenia, a region controlled by the Orontid Dynasty as one of the satrapies of the Achaemenid Empire. The Satrapy of Armenia was divided into two parts: the northern part and the southern part, with the cities of Erebuni (Yerevan) and Tushpa (Van) as their centres, respectively.
+
+Coins issued in 478 BC along with many other items found in the Erebuni Fortress, reveal the importance of Erebuni as a major centre for trade under Achaemenid rule.
+
+During the victorious period of Alexander the Great, and following the decline of the Achaemenid Empire, the Orontid rulers of the Armenian Satrapy achieved independence as a result of the Battle of Gaugamela in 331 BC, founding the Kingdom of Armenia. With the establishment of new cities such as Armavir, Zarehavan, Bagaran and Yervandashat, the importance of Erebuni had gradually declined.
+
+With the rise of the Artaxiad dynasty of Armenia who seized power in 189 BC, the Kingdom of Armenia greatly expanded to include major territories of Asia Minor, Atropatene, Iberia, Phoenicia and Syria. The Artaxiads considered Erebuni and Tushpa as cities of Persian heritage. Consequently, new cities and commercial centres were built by Kings Artaxias I, Artavasdes I and Tigranes the Great. Thus, with the dominance of cities such as Artaxata and Tigranocerta, Erebuni had significantly lost its importance as a central city.
+
+Under the rule of the Arsacid dynasty of Armenia (54–428 AD), many other cities around Erebuni including Vagharshapat and Dvin flourished. Consequently, Erebuni was completely neutralized, losing its role as an economic and strategic centre of Armenia. During the period of the Arsacid kings, Erebuni was only recorded in a Manichaean text of the 3rd century, where it is mentioned that one of the disciples of the prophet Mani founded a Manichaean community near the Christian community in Erebuni.
+
+According to Ashkharatsuyts, Erebuni was part of the Kotayk canton (Կոտայք գավառ, Kotayk gavar, not to be confused with the current Kotayk Province) of Ayrarat province, within Armenia Major.
+
+Armenia became a Christian nation in the early 4th century, during the reign of the Arsacid king Tiridates III.
+
+Following the partition of Armenia by the Byzantine and Sasanian empires in 387 and in 428, Erebuni and the entire territory of Eastern Armenia came under the rule of Sasanian Persia.[41] The Armenian territories formed the province of Persian Armenia within the Sasanian Empire.
+
+Due to the diminished role of Erebuni, as well as the absence of proper historical data, much of the city's history under the Sasanian rule is unknown.
+
+In 587 during the reign of emperor Maurice, Yerevan and much of Armenia came under Roman administration after the Romans defeated the Sassanid Persian Empire at the battle of the Blarathon. Soon after, Katoghike Tsiranavor Church in Avan was built between 595 and 602. Despite being partly damaged during the 1679 earthquake), it is the oldest surviving church within modern Yerevan city limits.
+
+The province of Persian Armenia (also known as Persarmenia) lasted until 646, when the province was dissolved with the Muslim conquest of Persia.
+
+In 658 AD, at the height of the Arab Islamic invasions, Erebuni-Yerevan was conquered during the Muslim conquest of Persia, as it was part of Persian-ruled Armenia. The city became part of the Emirate of Armenia under the Umayyad Caliphate. The city of Dvin was the centre of the newly created emirate. Starting from this period, as a result of the developing trade activities with the Arabs, the Armenian territories had gained strategic importance as a crossroads for the Arab caravan routes passing between Europe and India through the Arab-controlled Ararat plain of Armenia. Most probably, "Erebuni" has become known as "Yerevan" since at least the 7th century AD.
+
+After 2 centuries of Islamic rule over Armenia, the Bagratid prince Ashot I of Armenia led the revolution against the Abbasid Caliphate. Ashot I liberated Yerevan in 850, and was recognized as the Prince of Princes of Armenia by the Abbasid Caliph al-Musta'in in 862. Ashot was later crowned King of Armenia through the consent of Caliph al-Mu'tamid in 885. During the rule of the Bagratuni dynasty of Armenia between 885 and 1045, Yerevan was relatively a secure part of the Kingdom before falling to the Byzantines.
+
+However, Yerevan did not have any strategic role during the reign of the Bagratids, who developed many other cities of Ayrarat, such as Shirakavan, Dvin, and Ani.
+
+After a brief Byzantine rule over Armenia between 1045 and 1064, the invading Seljuks -led by Tughril and later by his successor Alp Arslan- ruled over the entire region, including Yerevan. However, with the establishment of the Zakarid Principality of Armenia in 1201 under the Georgian protectorate, the Armenian territories of Yerevan and Lori had significantly grown. After the Mongols captured Ani in 1236, Armenia turned into a Mongol protectorate as part of the Ilkhanate, and the Zakarids became vassals to the Mongols. After the fall of the Ilkhanate in the mid-14th century, the Zakarid princes ruled over Lori, Shirak and Ararat plain until 1360 when they fell to the invading Turkic tribes.
+
+During the last quarter of the 14th century, the Aq Qoyunlu Sunni Oghuz Turkic tribe took over Armenia, including Yerevan. In 1400, Timur invaded Armenia and Georgia, and captured more than 60,000 of the survived local people as slaves. Many districts including Yerevan were depopulated.[42]
+
+In 1410, Armenia fell under the control of the Kara Koyunlu Shia Oghuz Turkic tribe. According to the Armenian historian Thomas of Metsoph, although the Kara Koyunlu levied heavy taxes against the Armenians, the early years of their rule were relatively peaceful and some reconstruction of towns took place.[43] The Kara Koyunlus made Yerevan the centre of the newly formed Chukhur Saad administrative territory. The territory was named after a Turkic leader known as Emir Saad.
+
+However, this peaceful period was shattered with the rise of Qara Iskander between 1420 and 1436, who reportedly made Armenia a "desert" and subjected it to "devastation and plunder, to slaughter, and captivity".[44] The wars of Iskander and his eventual defeat against the Timurids, invited further destruction in Armenia, as many more Armenians were taken captive and sold into slavery and the land was subjected to outright pillaging, forcing many of them to leave the region.[45]
+
+Following the fall of the Armenian Kingdom of Cilicia in 1375, the seat of the Armenian Church was transferred from Sis back to Vagharshapat near Yerevan in 1441. Thus, Yerevan became the main economic, cultural and administrative centre in Armenia.
+
+In 1501–02, most of the Eastern Armenian territories including Yerevan were swiftly conquered by the emerging Safavid dynasty of Iran led by Shah Ismail I.[46] Soon after in 1502, Yerevan became the centre of the Erivan Beglarbegi, a new administrative territory of Iran formed by the Safavids. For the following 3 centuries, it remained, with brief intermissions, under the Iranian rule. Due to its strategic significance, Yerevan -known as Revan by the Ottomans- was initially often fought over, and passed back and forth, between the dominion of the rivaling Iranian and Ottoman Empire, until it permanently became controlled by the Safavids. In 1555, Iran had secured its legitimate possession over Yerevan with the Ottomans through the Treaty of Amasya.[47]
+
+In 1582–1583, the Ottomans led by Serdar Ferhad Pasha took brief control over Yerevan. Ferhad Pasha managed to build the Erivan Fortress on the ruins of one thousand-years old ancient Armenian fortress, on the shores of Hrazdan river.[48] However, Ottoman control ended in 1604 when the Persians regained Yerevan as a result of first Ottoman-Safavid War.
+
+Shah Abbas I of Persia who ruled between 1588 and 1629, ordered the deportation of hundreds of thousands of Armenians including citizens from Yerevan to mainland Persia. As a consequence, Yerevan significantly lost its Armenian population who had declined to 20%, while Muslims including Persians, Turks, Kurds and Tatars gained dominance with around 80% of the city's population. Muslims were either sedentary, semi-sedentary, or nomadic. Armenians mainly occupied the Kond neighbourhood of Yerevan and the rural suburbs around the city. However, the Armenians dominated over various professions and trade in the area and were of great economic significance to the Persian administration.[49]
+
+During the second Ottoman-Safavid War, Ottoman troops under the command of Sultan Murad IV conquered the city on 8 August 1635. Returning in triumph to Constantinople, he opened the "Yerevan Kiosk" (Revan Köşkü) in Topkapı Palace in 1636. However, Iranian troops under commanded by Shah Safi retook Yerevan on 1 April 1636. As a result of the Treaty of Zuhab in 1639, the Iranians reconfirmed their control over Eastern Armenia, including Yerevan. On 7 June 1679, a devastating earthquake razed the city to the ground.
+
+In 1724, the Erivan Fortress was besieged by the Ottoman army. After a period of resistance, the fortress fell to the Turks. As a result of the Ottoman invasion, the Erivan Beglarbegi of the Safavids was dissolved.
+
+Following a brief period of Ottoman rule over Eastern Armenia between 1724 and 1736, and as a result of the fall of the Safavid dynasty in 1736, Yerevan along with the adjacent territories became part of the newly formed administrative territory of Erivan Khanate under the Afsharid dynasty of Iran, which encompassed an area of 15,000 square kilometres (5,800 square miles). The Afsharids controlled Eastern Armenia from the mid 1730s until the 1790s. Following the fall of the Afsharids, the Qajar dynasty of Iran took control of Eastern Armenia until 1828, when the region was conquered by the Russian Empire after their victory over the Qajars that resulted in the Treaty of Turkmenchay of 1828.[50]
+
+During the second Russo-Persian War of the 19th century, the Russo-Persian War of 1826–28, Yerevan was captured by Russian troops under general Ivan Paskevich on 1 October 1827.[27][51][52] It was formally ceded by the Iranians in 1828, following the Treaty of Turkmenchay.[53] After 3 centuries of Iranian occupation, Yereven along with the rest of Eastern Armenia designated as the "Armenian Oblast", became part of the Russian Empire, a period that would last until the collapse of the Empire in 1917. The Russians sponsored the resettlement process of the Armenian population from Persia and Turkey. Due to the resettlement, the percentage of the Armenian population of Yerevan increased from 28% to 53.8%. The resettlement was intended to create Russian power bridgehead in the Middle East.[54] In 1829, Armenian repatriates from Persia were resettled in the city and a new quarter was built.
+
+Yerevan served as the seat of the newly formed Armenian Oblast between 1828 and 1840. By the time of Nicholas I's visit in 1837, Yerevan had become an uyezd. In 1840, the Armenian Oblast was dissolved and its territory incorporated into a new larger province; the Georgia-Imeretia Governorate. In 1850 the territory of the former oblast was reorganized into the Erivan Governorate, covering an area of 28,000 square kilometres (11,000 square miles). Yerevan was the centre of the newly established governorate.
+
+At that period, Yerevan was a small town with narrow roads and alleys, including the central quarter of Shahar, the Ghantar commercial centre, and the residential neighbourhoods of Kond, Dzoragyugh, Nork and Shentagh. During the 1840s and the 1850s, many schools were opened in the city. However, the first major plan of Yerevan was adopted in 1856, during which, Saint Hripsime and Saint Gayane women's colleges were founded and the English Park was opened. In 1863, the Astafyan Street was redeveloped and opened. In 1874, Zacharia Gevorkian opened Yerevan's first printing house, while the first theatre opened its doors in 1879.
+
+On 1 October 1879, Yerevan was granted the status of a city through a decree issued by Alexander II of Russia. In 1881, The Yerevan Teachers' Seminary and the Yerevan Brewery were opened, followed by the Tairyan's wine and brandy factory in 1887. Other factories for alcoholic beverages and mineral water were opened during the 1890s. The monumental church of Saint Gregory the Illuminator was opened in 1900. Electricity and telephone lines were introduced to the city in 1907 and 1913 respectively.
+
+In general, Yerevan had rapidly grown under the Russian rule, both economically and politically. Old buildings were torn down and new buildings of European style were erected instead.
+
+At the beginning of the 20th century, Yerevan city's population was over 29,000.[55] In 1902, a railway line linked Yerevan with Alexandropol, Tiflis and Julfa. In the same year, Yerevan's first public library was opened. In 1905, the grandnephew of Napoleon I; prince Louis Joseph Jérôme Napoléon (1864–1932) was appointed as governor of Yerevan province.[56] In 1913, for the first time in the city, a telephone line with eighty subscribers became operational.
+
+Yerevan served as the centre of the governorate until 1917, when Erivan governorate was dissolved with the collapse of the Russian Empire.
+
+At the beginning of the 20th century, Yerevan was a small city with a population of 30,000.[57] In 1917, the Russian Empire ended with the October Revolution. In the aftermath, Armenian, Georgian and Muslim leaders of Transcaucasia united to form the Transcaucasian Federation and proclaimed Transcaucasia's secession.
+
+The Federation, however, was short-lived. After gaining control over Alexandropol, the Turkish army was advancing towards the south and east to eliminate the center of Armenian resistance based in Yerevan. On 21 May 1918, the Turks started their campaign moving towards Yerevan via Sardarabad. Catholicos Gevorg V ordered that church bells peal for 6 days as Armenians from all walks of life – peasants, poets, blacksmiths, and even the clergymen – rallied to form organized military units.[58] Civilians, including children, aided in the effort as well, as "Carts drawn by oxen, water buffalo, and cows jammed the roads bringing food, provisions, ammunition, and volunteers from the vicinity" of Yerevan.[59]
+
+By the end of May 1918, Armenians were able to defeat the Turkish army in the battles of Sardarabad, Abaran and Karakilisa. Thus, on 28 May 1918, the Dashnak leader Aram Manukian declared the independence of Armenia. Subsequently, Yerevan became the capital and the center of the newly founded Republic of Armenia, although the members of the Armenian National Council were yet to stay in Tiflis until their arrival in Yerevan to form the government in the summer of the same year. Armenia became a parliamentary republic with four administrative divisions. The capital Yerevan was part of the Araratian Province. At the time, Yerevan received more than 75,000 refugees from Western Armenia, who escaped the massacres perpetrated by the Ottoman Turks during the Armenian Genocide.
+
+On 26 May 1919, the government passed a law to open the Yerevan State University, which was located on the main Astafyan (now Abovyan) street of Yerevan.
+
+After the signing of the Treaty of Sèvres in 1920, Armenia was granted formal international recognition. The United States, as well as many South American countries, officially opened diplomatic channels with the government of independent Armenia. Yerevan had also opened representatives in Great Britain, Italy, Germany, Serbia, Greece, Iran and Japan.
+
+However, after the short period of independence, Yerevan fell to the Bolsheviks, and Armenia was incorporated into Soviet Russia on 2 December 1920. Although nationalist forces managed to retake the city in February 1921 and successfully released all the imprisoned political and military figures, the city's nationalist elite were once again defeated by the Soviet forces on 2 April 1921.
+
+The Red Soviet Army invaded Armenia on 29 November 1920 from the northeast. On 2 December 1920, Yerevan along with the other territories of the Republic of Armenia, became part of Soviet Russia, known as the Armenian Soviet Socialist Republic. However, the Armenian SSR formed the Transcaucasian SFSR (TSFSR) together with the Georgian Soviet Socialist Republic and the Azerbaijan Soviet Socialist Republic, between 1922 and 1936.
+
+Under the Soviet rule, Yerevan became the first among the cities in the Soviet Union for which a general plan was developed. The "General Plan of Yerevan" developed by the academician Alexander Tamanian, was approved in 1924. It was initially designed for a population of 150,000. The city was quickly transformed into a modern industrial metropolis of over one million people. New educational, scientific and cultural institutions were founded as well.
+
+Tamanian incorporated national traditions with contemporary urban construction. His design presented a radial-circular arrangement that overlaid the existing city and incorporated much of its existing street plan. As a result, many historic buildings were demolished, including churches, mosques, the Persian fortress, baths, bazaars and caravanserais. Many of the districts around central Yerevan were named after former Armenian communities that were destroyed by the Ottoman Turks during the Armenian Genocide. The districts of Arabkir, Malatia-Sebastia and Nork Marash, for example, were named after the towns Arabkir, Malatya, Sebastia, and Marash, respectively. After the end of World War II, German POWs were used to help in the construction of new buildings and structures, such as the Kievyan Bridge.
+
+Within the years, the central Kentron district has become the most developed area in Yerevan, something that created a significant gap compared with other districts in the city. Most of the educational, cultural and scientific institutions were centred in the Kentron district.
+
+In 1965, during the commemorations of the fiftieth anniversary of the Armenian Genocide, Yerevan was the location of a demonstration, the first such demonstration in the Soviet Union, to demand recognition of the Genocide by the Soviet authorities.[60] In 1968, the city's 2,750th anniversary was commemorated.
+
+Yerevan played a key role in the Armenian national democratic movement that emerged during the Gorbachev era of the 1980s. The reforms of Glasnost and Perestroika opened questions on issues such as the status of Nagorno-Karabakh, the environment, Russification, corruption, democracy, and eventually independence. At the beginning of 1988, nearly one million Armenians from several regions of Armenia engaged in demonstrations concerning these subjects, centered in the city's Theater Square (currently Freedom Square).[61]
+
+Following the dissolution of the Soviet Union, Yerevan became the capital of Armenia on 21 September 1991.[62] Maintaining supplies of gas and electricity proved difficult; constant electricity was not restored until 1996 amidst the chaos of the badly instigated and planned transition to a market-based economy.
+
+Since 2000, central Yerevan has been transformed into a vast construction site, with cranes erected all over the Kentron district. Officially, the scores of multi-storied buildings are part of large-scale urban planning projects. Roughly $1.8 billion was spent on such construction in 2006, according to the national statistical service. Prices for downtown apartments have increased by about ten times during the first decade of the 21st century. Many new streets and avenues were opened, such as the Argishti street, Italy street, Saralanj Avenue, Monte Melkonian Avenue, and the Northern Avenue.
+
+However, as a result of this construction boom, the majority of the historic buildings located on the central Aram Street, were either entirely destroyed or transformed into modern residential buildings through the construction of additional floors. Only a few structures were preserved, mainly in the portion that extends between Abovyan Street and Mashtots Avenue.
+
+The first major post-independence protest in Yerevan took place in September 1996, after the announcement of incumbent Levon Ter-Petrosyan's victory in the presidential election. Major opposition parties of the time, consolidated around the former Karabakh Committee member and former Prime Minister Vazgen Manukyan, organized mass demonstrations between 23 and 25 September, claiming electoral fraud by Ter-Petrosyan.[63] An estimated of 200,000 people gathered in the Freedom Square to protest the election results.[64] After a series of riot and violent protests around the Parliament building on 25 September, the government sent tanks and troops to Yerevan to enforce the ban on rallies and demonstrations on the following day.[65] Prime Minister Vazgen Sargsyan and Minister of National Security Serzh Sargsyan announced on the Public Television of Armenia that their respective agencies have prevented an attempted coup d'état.[66]
+
+In February 2008, unrest in the capital between the authorities and opposition demonstrators led by ex-President Levon Ter-Petrosyan took place after the 2008 Armenian presidential election. The events resulted in 10 deaths[67] and a subsequent 20-day state of emergency declared by President Robert Kocharyan.[68]
+
+In July 2016, a group of armed men calling themselves the Daredevils of Sassoun (Armenian: Սասնա Ծռեր Sasna Tsrrer) stormed a police station in Erebuni District of Yerevan, taking several hostages, demanding the release of opposition leader Jirair Sefilian and the resignation of President Serzh Sargsyan. 3 policeman were killed as a result of the attack.[69] Many anti-government protestors held rallies in solidarity with the gunmen.[70] However, after 2 weeks of negotiations, the crisis ended and the gunmen surrendered.
+
+Yerevan has an average height of 990 m (3,248.03 ft), with a minimum of 865 m (2,837.93 ft) and a maximum of 1,390 m (4,560.37 ft) above sea level at the southwest and the northeast respectively.[71] It is located on to the edge of the Hrazdan River, northeast of the Ararat plain (Ararat Valley), to the center-west of the country. The upper part of the city is surrounded with mountains on three sides while it descends to the banks of the river Hrazdan at the south. The Hrazdan divides Yerevan into two parts through a picturesque canyon.
+
+Historically, the city is situated at the heart of the Armenian Highland,[72] in Kotayk canton (Armenian: Կոտայք գավառ Kotayk gavar, not to be confused with the current Kotayk Province) of Ayrarat province, within Armenia Major.
+
+As the capital of Armenia, Yerevan is not part of any marz ("province"). Instead, it is bordered with the following provinces: Kotayk from the north and the east, Ararat from the south and the south-west, Armavir from the west and Aragatsotn from the north-west.
+
+The Erebuni State Reserve formed in 1981, is located around 8 km southeast of the city centre within the Erebuni District of the city. At a height between 1300 and 1450 meters above sea level, the reserve occupies an area of 120 hectares, mainly consisted of semi-deserted mountains-steppe.[73]
+
+Yerevan features a continental influenced steppe climate (Köppen climate classification: BSk or "cold semi-arid climate"), with long, hot, dry summers and short, but cold and snowy winters. This is attributed to Yerevan being on a plain surrounded by mountains and to its distance from the sea and its effects. The summers are usually very hot with the temperature in August reaching up to 40 °C (104 °F), and winters generally carry snowfall and freezing temperatures with January often being as cold as −15 °C (5 °F) and lower. The amount of precipitation is small, amounting annually to about 318 millimetres (12.5 in). Yerevan experiences an average of 2,700 sunlight hours per year.[71]
+
+The Yerevan TV Tower is the tallest structure in the city, and one of the tallest structures in the Transcaucasian region.
+
+The Republic Square, the Yerevan Opera Theatre, and the Yerevan Cascade are among the main landmarks at the centre of Yerevan, mainly developed based on the original design of the academician Alexander Tamanian, and the revised plan of architect Jim Torosyan.
+
+A major redevelopment process has been launched in Yerevan since 2000. As a result, many historic structures have been demolished and replaced with new buildings. This urban renewal plan has been met with opposition[76] and criticism from some residents, as the projects destroy historic buildings dating back to the period of the Russian Empire, and often leave residents homeless.[77][78][79] Downtown houses deemed too small are increasingly demolished and replaced by high-rise buildings.
+
+The Saint Gregory Cathedral, the new building of Yerevan City Council, the new section of Matenadaran institute, the new terminal of Zvartnots International Airport, the Cafesjian Center for the Arts at the Cascade, the Northern Avenue, and the new government complex of ministries are among the major construction projects fulfilled during the first two decades of the 21st century.
+
+Aram Street of old Yerevan and the newly built Northern Avenue are respectively among the notable examples featuring the traditional and modern architectural characteristics of Yerevan.
+
+As of May 2017, Yerevan is home to 4,883 residential apartment buildings, and 65,199 street lamps installed on 39,799 street light posts, covering a total length of 1,514 km. The city has 1,080 streets with a total length of 750 km.[80]
+
+Yerevan is a densely built city but still offers several public parks throughout its districts, graced with mid-sized green gardens. The public park of Erebuni District along with its artificial lake is the oldest garden in the city. Occupying an area of 17 hectares, the origins of the park and the artificial lake date back to the period of king Argishti I of Urartu during the 8th century BC. In 2011, the garden was entirely remodeled and named as Lyon Park, to become a symbol of the partnership between the cities of Lyon and Yerevan.[81]
+
+The Lovers' Park on Marshal Baghramyan Avenue and the English Park at the centre of the city, dating back to the 18th and 19th centuries respectively, are among the most popular parks for the Yerevantsis. The Yerevan Botanical Garden opened in 1935, the Victory park formed in the 1950s and the Circular Park are among the largest green spaces of the city.
+
+Formed in the 1960s, the Yerevan Opera Theatre park along with its artificial Swan Lake is also among the favorite green spaces of the city. In 2019 some of the public space of the park leased to restaurants was reclaimed allowing for improved landscape design.[82] A public ice-skating arena is operated in the park's lake area during winters.
+
+The Yerevan Lake is an artificial reservoir opened in 1967 on Hrazdan riverbed at the south of the city centre, with a surface of 0.65 km2 (0.25 sq mi).
+
+Each administrative district of Yerevan has its own public park, such as the Buenos Aires Park and Tumanyan Park in Ajapnyak, Komitas park in Shengavit, Vahan Zatikian park in Malatia-Sebastia, David Anhaght park in Kanaker-Zeytun, the Family park in Avan, and Fridtjof Nansen park in Nor Nork.
+
+Yerevan has been the capital of Armenia since the independence of the First Republic in 1918. Situated in the Ararat plain, the historic lands of Armenia, it served as the best logical choice for capital of the young republic at the time.
+
+When Armenia became a republic of the Soviet Union, Yerevan remained as capital and accommodated all the political and diplomatic institutions in the republic. In 1991 with the independence of Armenia, Yerevan continued with its status as the political and cultural centre of the country, being home to all the national institutions: the Government House, the National Assembly, the Presidential Palace, the Central Bank, the Constitutional Court, all ministries, judicial bodies and other government organizations.
+
+Yerevan received the status of a city on 1 October 1879, upon a decree issued by Tsar Alexander II of Russia. The first city council formed was headed by Hovhannes Ghorghanyan, who became the first mayor of Yerevan.
+
+The Constitution of the Republic of Armenia adopted on 5 July 1995, granted Yerevan the status of a marz (մարզ, province).[83] Therefore, Yerevan functions similarly to the provinces of Armenia with a few specifications.[84]
+The administrative authority of Yerevan is thus represented by:
+
+In the modified Constitution of 27 November 2005, Yerevan city was turned into a "community" (համայնք, hamaynk); since, the Constitution declares that this community has to be led by a mayor, elected directly or indirectly, and that the city needs to be governed by a specific law.[88] The first election of the Yerevan City Council took place in 2009 and won by the ruling Republican Party of Armenia.[89][90]
+
+In addition to the national police and road police, Yerevan has its own municipal police. All three bodies cooperate to maintain law in the city.
+
+Yerevan is divided into twelve "administrative districts" (վարչական շրջան, varčakan šrĵan)[91] each with an elected leader. The total area of the 12 districts of Yerevan is 223 square kilometres (86 square miles).[92][93][94]
+
+Originally a small town, Yerevan became the capital of Armenia and a large city with over one million inhabitants. Until the fall of the Soviet Union, the majority of the population of Yerevan were Armenians with minorities of Russians, Kurds, Azerbaijanis and Iranians present as well. However, with the breakout of the Nagorno-Karabakh War from 1988 to 1994, the Azerbaijani minority diminished in the country in what was part of population exchanges between Armenia and Azerbaijan. A big part of the Russian minority also fled the country during the 1990s economic crisis in the country. Today, the population of Yerevan is overwhelmingly Armenian.
+
+After the collapse of the Soviet Union, due to economic crises, thousands fled Armenia, mostly to Russia, North America and Europe. The population of Yerevan fell from 1,250,000 in 1989[71] to 1,103,488 in 2001[107] and to 1,091,235 in 2003.[108] However, the population of Yerevan has been increasing since. In 2007, the capital had 1,107,800 inhabitants.
+
+Yerevantsis in general use the Yerevan dialect, an Eastern Armenian dialect most probably formed during the 13th century. It is currently spoken in and around Yerevan, including the towns of Vagharshapat and Ashtarak. Classical Armenian (Grabar) words compose a significant part of the dialect's vocabulary.[109] Throughout the history, it was influenced by several languages, especially Russian and Persian and loan words have significant presence in it today. It is currently the most widespread Armenian dialect.[110]
+
+Yerevan was inhabited first by Armenians and remained homogeneous until the 15th century.[95][96][111][better source needed] The population of the Erivan Fortress, founded in the 1580s, was mainly composed of Muslim soldiers, estimated two to three thousand.[95] The city itself was mainly populated by Armenians. French traveler Jean-Baptiste Tavernier, who visited Yerevan possibly up to six times between 1631 and 1668, states that the city is exclusively populated by Armenians.[112][better source needed] During the 1720s Ottoman–Persian War[clarification needed] its absolute majority were Armenians.[96] The demographics of the region changed because of a series of wars between the Ottoman Empire, Iran and Russia. By the early 19th century, Yerevan had a Muslim majority.
+
+Until the Sovietization of Armenia, Yerevan was a multicultural city, mainly with an Armenian and "Caucasian Tatar" (modern-day Azerbaijani) population. After the Armenian Genocide, many refugees from what Armenians call Western Armenia (nowadays Turkey, then Ottoman Empire) escaped to Eastern Armenia. In 1919, about 75,000 Armenian refugees from the Ottoman Empire arrived in Yerevan, mostly from the Vaspurakan region (city of Van and surroundings). A significant part of these refugees died of typhus and other diseases.[113]
+
+From 1921 to 1936, about 42,000 ethnic Armenians from Iraq, Turkey, Iran, Greece, Syria, France, Bulgaria etc. went to Soviet Armenia, with most of them settling in Yerevan. The second wave of repatriation occurred from 1946 to 1948, when about 100,000 ethnic Armenians from Iran, Syria, Lebanon, Greece, Bulgaria, Romania, Cyprus, Palestine, Iraq, Egypt, France, United States etc. moved to Soviet Armenia, again most of whom settled in Yerevan. Thus, the ethnic makeup of Yerevan became more monoethnic during the first 3 decades in the Soviet Union. In the late 1980s and the early 1990s, the remaining 2,000 Azeris left the city, because of the Nagorno-Karabakh conflict.
+
+Armenian Apostolic Christianity is the predominant religion in Armenia. The 5th-century Saint Paul and Peter Church demolished in November 1930 by the Soviets, was among the earliest churches ever built in Erebuni-Yerevan. Many of the ancient Armenian and medieval churches of the city were destroyed by the Soviets in the 1930s during the Great Purge.
+
+The regulating body of the Armenian Church in Yerevan is the Araratian Pontifical Diocese, with the Surp Sarkis Cathedral being the seat of the diocese. It is the largest diocese of the Armenian Church and one of the oldest dioceses in the world, covering the city of Yerevan and the Ararat Province of Armenia.[19]
+
+Yerevan is currently home to the largest Armenian church in the world, the Cathedral of Saint Gregory the Illuminator. It was consecrated in 2001, during the 1700th anniversary of the establishment of the Armenian Church and the adoption of Christianity as the national religion in Armenia.
+
+As of 2017, Yerevan has 17 active Armenian churches as well as four chapels.
+
+After the capture of Yerevan by the Russians as a result of the Russo-Persian War of 1826–28, many Russian Orthodox churches were built in the city under the orders of the Russian commander General Ivan Paskevich. The Saint Nikolai Cathedral opened during the second half of the 19th century, was the largest Russian church in the city. The Church of the Intercession of the Holy Mother of God was opened in 1916 in Kanaker-Zeytun.[114]
+
+However, most of the churches were either closed or demolished by the Soviets during the 1930s. The Saint Nikolai Cathedral was entirely destroyed in 1931, while the Church of the Intercession of the Holy Mother of God was closed and converted first into a warehouse and later into a club for the military personnel. Religious services resumed in the church it in 1991, and in 2004 a cupola and a belfry were added to the building.[115] In 2010, the groundbreaking ceremony of the new Holy Cross Russian Orthodox church took place with the presence of Patriarch Kirill I of Moscow. The church was eventually consecrated on 7 October 2017, with the presence of Catholicos Karekin II, Russian bishops and the church benefactor Ara Abramyan.
+
+According to Ivan Chopin, there were eight mosques in Yerevan in the middle of the 19th century.[116][117] The 18th-century Blue Mosque of Yerevan was restored and reopened in the 1990s, with Iranian funding,[118] and is currently the only active mosque in Armenia, mainly serving the Iranian Shia visitors.
+
+Yerevan is home to tiny Yezidi, Molokan, Neopagan, Bahai and Jewish communities, with the Jewish community being represented by the Jewish Council of Armenia. A variety of nontrinitarian communities, considered dangerous sects by the Armenian Apostolic Church,[119] are also found in the city, including Jehovah's Witnesses, Mormons, Seventh-day Adventists and Word of Life.[120]
+
+Medical services in Armenia – except from maternity – are not subsidized by the government. However, the government annually allocates a certain amount from the state budget for the medical needs of the socially vulnerable groups.
+
+Yerevan is a major healthcare and medical service centre in the region. Several hospitals of Yerevan refurbished with modern technologies, provide healthcare and medical researches, such as Shengavit Medical Center, Erebouni Medical Center, Izmirlian Medical Center, Saint Gregory the Illuminator Medical Center, Nork-Marash Medical Center, Armenia Republican Medical Center, Astghik Medical Centre, Armenian American Wellness Center, and Mkhitar Heratsi Hospital Complex of the Yerevan State Medical University. The municipality runs 39 polyclinics/medical centres throughout the city.
+
+The Research Center of Maternal and Child Health Protection is operating in Yerevan since 1937, while the Armenicum Clinical Center was opened in 1999,[121] where researches are conducted mainly about infectious diseases and associated researches, including HIV, immunodeficiency and hepatitis.
+
+The Liqvor Pharmaceuticals Factory operating since 1991 in Yerevan, is currently the largest medicines manufacturer of Armenia.[122]
+
+Yerevan is Armenia's principal cultural, artistic, and industrial center, as well as the seat of the national government with a large number of museums, important monuments and the national public library. It also hosts Vardavar the most widely celebrated festival among Armenians and is one of the historic centres of traditional Armenian carpet.
+
+Yerevan is home to a large number of museums, art galleries and libraries. The most prominent of these are the National Gallery of Armenia, the History Museum of Armenia, the Cafesjian Museum of Art, the Matenadaran library of ancient manuscripts, and the Armenian Genocide museum of Tsitsernakaberd complex.
+
+Founded in 1921, the National Gallery of Armenia and the History Museum of Armenia are the principal museums of the city. In addition to having a permanent exposition of works of Armenian painters, the gallery houses a collection of paintings, drawings and sculptures issued from German, American, Austrian, Belgian, Spanish, French, Hungarian, Italian, Dutch, Russian and Swiss artists.[123] It usually hosts temporary expositions.
+
+The Armenian Genocide museum is found at the foot of Tsitsernakaberd memorial and features numerous eyewitness accounts, texts and photographs from the time. It comprises a memorial stone made of three parts, the latter of which is dedicated to the intellectual and political figures who, as the museum's site says, "raised their protest against the Genocide committed against the Armenians by the Turks. Among them there are Armin T. Wegner, Hedvig Büll, Henry Morgenthau Sr., Franz Werfel, Johannes Lepsius, James Bryce, Anatole France, Giacomo Gorrini, Benedict XV, Fritjof Nansen, and others.
+
+Cafesjian Museum of Art within the Cascade complex, is an art centre opened on 7 November 2009. It showcases a massive collection glass artwork, particularly the works of the Czech artists Stanislav Libenský and Jaroslava Brychtová. The front gardens showcase sculptures from Gerard L. Cafesjian's collection.
+
+The Erebuni Museum founded in 1968, is an archaeological museum housing Urartian artifacts found during excavations at the Erebuni Fortress. The Yerevan History Museum and the Armenian Revolutionary Federation History Museum are among the prominent museums that feature the history of Yerevan and the First Republic of Armenia respectively. The Military Museum within the Mother Armenia complex is about the participation of Armenian soldiers in World War II and Nagorno-Karabakh War.
+
+The city is also home to a large number of art museums. Sergei Parajanov Museum opened in 1988 is dedicated to Sergei Parajanov's art works in cinema and painting.[124] Komitas Museum opened in 2015, is a musical art museum devoted to the renowned Armenian composer Komitas. Charents Museum of Literature and Arts opened in 1921, Modern Art Museum of Yerevan opened in 1972, and the Middle East Art Museum opened in 1993, are also among the notable arte museums of the city.[125]
+
+Biographical museums are also common in Yerevan. Many renowned Armenian poets, painters and musicians are honored with house-museums in their memory, such as poet Hovhannes Tumanyan, composer Aram Khachaturian, painter Martiros Saryan, novelist Khachatur Abovian, and French-Armenian singer Charles Aznavour.
+
+Recently, many museums of science and technology have opened in Yerevan, such as the Museum of Armenian Medicine (1999), the Space Museum of Yerevan (2001), Museum of Science and Technology (2008), Museum of Communications (2012) and the Little Einstein Interactive Science Museum (2016).
+
+The National Library of Armenia located on Teryan Street of Yerevan, is the public library of the city and the entire republic. It was founded in 1832 and is operating in its current building since 1939. Another national library of Yerevan is the Khnko Aper Children's Library, founded in 1933. Other major public libraries include the Avetik Isahakyan Central Library founded in 1935, the Republican Library of Medical Sciences founded in 1939, the Library of Science and Technology founded in 1957, and the Musical Library founded in 1965. In addition, each administrative district of Yerevan has its own public library (usually more than one library).
+
+The Matenadaran is a library-museum and a research centre, regrouping 17,000 ancient manuscripts and several bibles from the Middle Ages. Its archives hold a rich collection of valuable ancient Armenian, Ancient Greek, Aramaic, Assyrian, Hebrew, Latin, Middle and Modern Persian manuscripts. It is located on Mashtots Avenue at central Yerevan.
+
+On 6 June 2010, Yerevan was named as the 2012 World Book Capital by the United Nations Educational, Scientific and Cultural Organization (UNESCO). The Armenian capital was chosen for the quality and variety of the programme it presented to the selection committee, which met at UNESCO's headquarters in Paris on 2 July 2010.
+
+The National Archives of Armenia founded in 1923, is a scientific research centre and depositary, with a collection of around 3.5 million units of valuable documents.
+
+Yerevan is one of the historic centres of traditional Armenian carpet. Various rug fragments have been excavated in areas around Yerevan dating back to the 7th century BC or earlier. The tradition was further developed from the 16th century when Yerevan became the central city of Persian Armenia. However, carpet manufacturing in the city was greatly enriched with the flock of Western Armenian migrants from the Ottoman Empire throughout the 19th century, and the arrival of Armenian refugees escaping the genocide in the early 20th century. Currently, the city is home to the Arm Carpet factory opened in 1924, as well as the Tufenkian handmade carpets (since 1994), and Megerian handmade carpets (since 2000).
+
+The Yerevan Vernissage open-air exhibition-market formed in the late 1980s on Aram Street, features a large collection of different types of traditional Armenian hand-made art works, especially woodwork sculptures, rugs and carpets. On the other hand, the Saryan park located near the opera house, is famous for being a permanent venue where artists exhibit their paintings.
+
+The Armenian Center for Contemporary Experimental Art founded in 1992 in Yerevan,[126] is a creativity centre helping to exchange experience between professional artists in an appropriate atmosphere.[127]
+
+Jazz, classical, folk and traditional music are among several genres that are popular in the city of Yerevan. A large number of ensembles, orchestras and choirs of different types of Armenian and international music are active in the city.
+
+The Armenian Philharmonic Orchestra founded in 1925, is one of the oldest musical groups in Yerevan and modern Armenia. The Armenian National Radio Chamber Choir founded in 1929, won the First Prize of the Soviet Union in the 1931 competition of choirs among the republics of the Soviet Union. Folk and classical music of Armenia was taught in state-sponsored conservatoires during the Soviet days. The Sayat-Nova Armenian Folk Song Ensemble was founded in Yerevan in 1938. Currently directed by Tovmas Poghosyan, the ensemble performs the works of prominent Armenian gusans such as Sayat-Nova, Jivani, and Sheram.
+
+In 1939, the Armenian National Academic Theatre of Opera and Ballet was opened. It is home to the Aram Khatchaturian concert hall and the Alexander Spendiarian auditorium of the National Theatre of Opera and Ballet.
+
+The Komitas Chamber Music House opened in 1977, is the home of chamber music performers and lovers in Armenia. In 1983, the Karen Demirchyan Sports and Concerts Complex was opened. It is currently the largest indoor venue in Armenia.
+
+The National Chamber Orchestra of Armenia (founded in 1961), Yerevan State Brass Band (1964), Folk Instruments Orchestra of Armenia (1977), Gusan and Folk Song Ensemble of Armenia (1983), Hover Chamber Choir (1992), Shoghaken Folk Ensemble (1995), Yerevan State Chamber Choir (1996), State Orchestra of Armenian National Instruments (2004), and the Youth State Orchestra of Armenia (2005), are also among the famous musical ensembles of the city of Yerevan. The Ars lunga piano-cello duo achieved international fame since its foundation in 2009 in Yerevan.
+
+Armenian religious music remained liturgical until Komitas introduced polyphony by the end of the 19th century. Starting from the late 1950s, religious music became widely spread when Armenian chants (also known as sharakans) were performed by the soprano Lusine Zakaryan. The state-run Tagharan Ensemble of Yerevan founded in 1981 and currently directed by Sedrak Yerkanian, also performs ritual and ancient Armenian music.
+
+Jazz is also among the popular genres in Yerevan. The first jazz band in Yerevan was founded in 1936. Currently, many jazz and ethno jazz bands are active in Yerevan such as Time Report, Art Voices, and Nuance Jazz Band. The Malkhas jazz club founded by renowned artist Levon Malkhasian, is among the most popular clubs in the city. The[Yerevan Jazz Fest is an annual jazz festival taking place every autumn since 2015, organized by the Armenian Jazz Association with the support of the Yerevan Municipality.[128]
+
+Armenian rock has been originated in Yerevan in the mid 1960s, mainly through Arthur Meschian and his band Arakyalner (Disciples). In the early 1970s, there were a range of professional bands in Yerevan strong enough to compete with their Soviet counterparts. In post-Soviet Armenia, an Armenian progressive rock scene has been developed in Yerevan, mainly through Vahan Artsruni, the Oaksenham rock band, and the Dorians band. The Armenian Navy Band founded by Arto Tunçboyacıyan in 1998 is also famous for jazz, avant-garde and folk music. Reggae is also becoming popular in Yerevan mainly through the Reincarnation musical band.
+
+The Cafesjian Center for the Arts is known for its regularly programmed events including the "Cafesjian Classical Music Series" on the first Wednesday of each month, and the "Music Cascade" series of jazz, pop and rock music live concerts performed every Friday and Saturday.
+
+Open-air concerts are frequently held in curtain location in Yerevan during summer, such as the Cafesjian Sculpture Garden on Tamanyan Street, the Freedom Square near the Opera House, the Republic Square, etc. The famous KOHAR Symphony Orchestra and Choir occasionally performs open-air concerts in the city.
+
+Traditional dancing is very popular among Armenians. During the cool summertime of the Yerevan city, it is very common to find people dancing in groups at the Northern Avenue or the Tamanyan Street near the cascade.
+
+Professional dance groups were formed in Yerevan during the Soviet days. The first group was the Armenian Folk Music and Dance Ensemble founded in 1938 by Tatul Altunyan. It was followed by the State Dance Ensemble of Armenia in 1958. In 1963, the Berd Dance Ensemble was formed. The Barekamutyun State Dance Ensemble of Armenia was founded in 1987 by Norayr Mehrabyan.
+
+The Karin Traditional Song and Dance Ensemble founded in 2001 by Gagik Ginosyan is known for revitalizing and performing the ancient Armenian dances of the historical regions of the Armenian Highlands,[129] such as Hamshen, Mush, Sasun, Karin, etc.
+
+Yerevan is home to many theatre groups, mainly operating under the support of the ministry of culture. Theatre halls in the city organize several shows and performances throughout the year. Most prominent state-run theatres of Yerevan are the Sundukyan State Academic Theatre, Paronyan Musical Comedy Theatre, Stanislavski Russian Theatre, Hrachya Ghaplanyan Drama Theatre, and the Sos Sargsyan Hamazgayin State Theatre. The Edgar Elbakyan Theatre of Drama and Comedy is among the prominent theatres run by the private sector.
+
+Yerevan is also home to several specialized theatres such as the Tumanyan Puppet Theatre, Yerevan State Pantomime Theatre, and the Yerevan State Marionettes Theatre.
+
+Cinema in Armenia was born on 16 April 1923, when the Armenian State Committee of Cinema was established upon a decree issued by the Soviet Armenian government.
+
+In March 1924, the first Armenian film studio; Armenfilm (Armenian: Հայֆիլմ "Hayfilm," Russian: Арменкино "Armenkino") was opened in Yerevan, starting with a documentary film called Soviet Armenia. Namus was the first Armenian silent black and white film, directed by Hamo Beknazarian in 1925, based on a play of Alexander Shirvanzade, describing the ill fate of two lovers, who were engaged by their families to each other since childhood, but because of violations of namus (a tradition of honor), the girl was married by her father to another person. The first produced sound film was Pepo  directed by Hamo Beknazarian in 1935.
+
+Nowadays, Yerevan is home to many movie theatres including the Moscow Cinema, Nairi Cinema, Hayastan Cinema, Cinema Star multiplex cinemas of the Dalma Garden Mall, and the KinoPark multiplex cinemas of Yerevan Mall. Since 2004, the Moscow Cinema hosts the Golden Apricot Yerevan International Film Festival annually. The ReAnimania International Animation Film & Comics Art Festival of Yerevan launched in 2005, is also among the popular annual events in the city.[130]
+
+In addition to the art festivals, the city organizes many public celebrations that greatly attract the locals as well as the visitors. Vardavar is the most widely celebrated festival among Armenians, having it roots back to the pagan history of Armenia. It is celebrated 98 days (14 weeks) after Easter. During the day of Vardavar, people from a wide array of ages are allowed to douse strangers with water. It is common to see people pouring buckets of water from balconies on unsuspecting people walking below them. The Swan Lake of the Yerevan Opera is the most popular venue for the Vardavar celebrations.
+
+In August 2015, Teryan Cultural Centre supported by the Yerevan Municipality has launched its 1st Armenian traditional clothing festival known as the Yerevan Taraz Fest.[131]
+
+As one of the ancient winemaking regions, many wine festivals are celebrated in Armenia. Yerevan launched its 1st annual wine festivals known as the Yerevan Wine Days in May 2016.[132] The Watermelon Fest launched in 2013 is also becoming a popular event in the city. The Yerevan Beer Fest is held annually during the month of August. It was first organized in 2014.[133]
+
+Many public and private TV and radio channels operate in Yerevan. The Public TV of Armenia has been in service since 1956. It became a satellite television in 1996. Other satellite TVs include the Armenia TV owned by the Pan-Armenian Media Group, Kentron TV owned by Gagik Tsarukyan, Shant TV and Shant TV premium. On the other hand, Yerkir Media, Armenia 2, Shoghakat TV, Yerevan TV, 21TV and the TV channels of the Pan-Armenian Media Group are among the most notable local televisions of Yerevan.
+
+Notable newspapers published in Yerevan include the daily newspapers of Aravot, Azg, Golos Armenii and Hayastani Hanrapetutyun.
+
+Many of the structures of Yerevan had been destroyed either during foreign invasions or as a result of the devastating earthquake in 1679. However, some structures have remained moderately intact and were renovated during the following years.
+
+Erebuni Fortress, also known as Arin Berd, is the hill where the city of Yerevan was founded in 782 BC by King Argishti I. The remains of other structures from earlier periods are also found in Shengavit.
+
+The 4th-century chapel of the Holy Mother of God and the 6th-century Tsiranavor Church both located in Avan District at the north of Yerevan, are among the oldest surviving Christian structures of the city. Originally a suburb at the north of Yerevan, Avan was eventually absorbed by the city's gradual expansion. The district is also home to the remains of Surp Hovhannes Chapel dating back to the 12–13th centuries.
+
+Katoghike Church; a medieval chapel in the centre of Yerevan built in 1264, is one of the best preserved churches of the city.[134] Zoravor Surp Astvatsatsin Church is also among the best surviving churches of Yerevan, built 1693–94 right after the devastating earthquake, on the ruins of a medieval church. Saint Sarkis Cathedral rebuilt in 1835–42, is the seat of Araratian Pontifical Diocese of the Armenian Church.
+
+The Blue Mosque or "Gök Jami", built between 1764 and 1768 at the centre of the city, is currently the only operating mosque in Armenia.
+
+The Red Bridge of Hrazdan River is a 17th-century structure, built after the 1679 earthquake and later reconstructed in 1830.
+
+Yerevan Opera Theater or the Armenian National Academic Opera and Ballet Theatre opened in 1933, is a major landmark in the city along with the Mesrop Mashtots Matenadaran opened in 1959, and Tsitsernakaberd monument of the Armenian Genocide opened in 1967.
+
+Moscow Cinema, opened in 1937 on the site of Saint Paul and Peter Church of the 5th century, is an important example of the Soviet-era architecture. In 1959, a monument was erected near the Yerevan Railway Station dedicated to the legendary Armenian hero David of Sassoun. The monumental statue of Mother Armenia is a female personification of the Armenian nation, erected in 1967, replacing the huge statue of Joseph Stalin in the Victory park.
+
+Komitas Pantheon is a cemetery opened in 1936 where many famous Armenians are buried, while the Yerablur Pantheon, is a military cemetery where over 1,000 Armenian martyrs of the Nagorno-Karabakh War are buried since 1990.
+
+Many new notable buildings were constructed after the independence of Armenia such as the Yerevan Cascade, and the Saint Gregory Cathedral opened in 2001 to commemorate the 1700th anniversary of Christianity in Armenia. In May 2016, a monumental statue of the prominent Armenian statesman and military leader Garegin Nzhdeh was erected at the centre of Yerevan.
+
+Yerevan is served by the Zvartnots International Airport, located 12 kilometres (7 miles) west of the city center. It is the primary airport of the country. Inaugurated in 1961 during the Soviet era, Zvartnots airport was renovated for the first time in 1985 and a second time in 2002 in order to adapt to international norms. It went through a facelift starting in 2004 with the construction of a new terminal. The first phase of the construction ended in September 2006 with the opening of the arrivals zone. A second section designated for departures was inaugurated in May 2007. The departure terminal is anticipated, October 2011 housing state of the art facilities and technology. This will make Yerevan Zvartnots International Airport, the largest, busiest and most modern airport in the entire Caucasus. Currently there are no national airlines operating in Armenia.[135] The entire project costs more than US$100 million.
+
+A second airport, Erebuni Airport, is located just south of the city. Since the independence, "Erebuni" is mainly used for military or private flights. The Armenian Air Force has equally installed its base there and there are several MiG-29s stationed on Erebuni's tarmac.
+
+Public transport in Yerevan is heavily privatized and mostly handled by around 60 private operators. As of May 2017, 39 city bus lines are being operated throughout Yerevan.[136] These lines mostly consist of about 425 Bogdan, Higer City Bus and Hyundai County buses. However, the market share these buses in public transit is only about 39.1%.
+
+But the 50.4% of public transit is still served by "public vans", locally known as marshrutka. These are about 1210 Russian-made GAZelle vans with 13 seats, that operate same way as buses, having 79 different lines with certain routes and same stops. According to Yerevan Municipality office, in future, marshrutkas should be replaced by ordinary larger buses. Despite having about 13 seats, the limit of passengers is not controlled, so usually these vans carry many more people who stand inside.
+
+The Yerevan trolleybus system has been operating since 1949. Some old Soviet-era trolleybuses have been replaced with comparably new ones. As of May 2017, only 5 trolleybus lines are in operation (2.6% share), with around 45 units in service. The trolleybus system is owned and operated by the municipality.
+
+The tram network that operated in Yerevan since 1906 was decommissioned in January 2004. Its operation had a cost 2.4 times higher than the generated profits, which pushed the municipality to shut down the network,[137] despite a last-ditch effort to save it towards the end of 2003. Since the closure, the rails have been dismantled and sold.
+
+Due to being dispersed among dozens of private operators, the transportation is barely regulated, with only trip fee is being a subject of regulation. Thus, the quality of vehicles is often inadequate, with no certain regulations for safety. Unlike the majority of world capitals, there is no established ticketing system in Yerevan's public transportation. Passengers need to pay the money directly to the driver when getting out of the vehicle. The fare -being one of the few things that is regulated- is fixed and controlled by authorities. A one-way trip costs AMD 100 (around US$0.21) for all buses and public vans, while it is AMD 50 for trolleybuses.
+
+The central station in Nor Kilikia neighborhood serves as bus terminal for inter-city transport, serving outbound routes towards practically all the cities of Armenia as well as abroad, notably Tbilisi and Tabriz.
+
+The Yerevan Metro named after Karen Demirchyan, (Armenian: Կարեն Դեմիրճյանի անվան Երեւանի մետրոպոլիտեն կայարան (Karen Dyemirchyani anvan Yerevani metropoliten kayaran)) is a rapid transit system that serves the capital city since 1981. It has a single line of 12.1 km (7.5 mi) length with 10 active stations and 45 units in service. The interiors of the stations resemble that of the former western Soviet nations, with chandeliers hanging from the corridors. The metro stations had most of their names changed after the collapse of the Soviet Union and the independence of the Republic of Armenia.
+
+A northeastern extension of the line with two new stations is currently being developed. The construction of the first station (Ajapnyak) and of the one-kilometre (0.62-mile) tunnel linking it to the rest of the network will cost US$18 million.[138] The time of the end of the project has not yet been defined. Another long-term project is the construction of two new lines, but these have been suspended due to lack of finance.
+
+The system transports more than 60,000 people on a daily basis.
+
+Yerevan has a single central railway station (several railway stations of suburbs have not been used since 1990) that is connected to the metro via the Sasuntsi Davit station. The railway station is made in Soviet-style architecture with its long point on the building roof, representing the symbols of communism: red star, hammer and sickle. Due to the Turkish and Azerbaijani blockades of Armenia, there is only one international train that passes by once every two days, with neighboring Georgia being its destination. For a sum of 9 000 to 18 000 dram, it is possible to take the night train to the Georgian capital, Tbilisi.[139] This train then continues to its destination of Batumi, on the shores of the Black sea in the summer season.
+
+The only railway that goes to Iran to the south passes by the closed border of Nakhichevan. For this reason, there are no trains that go south from Yerevan. A construction project on a new railway line connecting Armenia and Iran directly is currently being studied.
+
+During the first decade of the 21st century, the South Caucasus Railway CJSC—which is the current operator of the railway system in Armenia—announced its readiness to put the Yerevan-Gyumri-Kars railway line in service in case the Armenian-Turkish protocols are ratified and the opening of the borders between the two countries is achieved.
+
+As of July 2017, the following railway trips are scheduled from and to Yerevan:
+
+The Yerevan-Ararat-Yerevan route is temporarily not in operation, while the Yerevan-Tbilisi-Yerevan route will operate starting from 2 October 2017.
+
+Armenia is among the top 10 safest countries where one can wander around and go home alone safely at night. Yerevan prides itself on having connections 24/7 as taxis are available at any time of the day or night.[144] Taxicab service companies cover the entire city in addition to many online taxi service providers, including GG Taxi, Utaxi and Yandex.Taxi.
+
+As of 2013[update], the share of Yerevan in the annual total industrial product of Armenia is 41%.[145] The industry of Yerevan is quite diversified including chemicals, primary metals and steel products, machinery, rubber products, plastics, rugs and carpets, textiles, clothing and footwear, jewellery, wood products and furniture, building materials and stone-processing, alcoholic beverages, mineral water, dairy product and processed food. Even though the economic crisis of the '90s ravaged the industry of the country, several factories remain always in service, notably in the petrochemical and the aluminium sectors.
+
+Armenian beverages, especially Armenian cognac and beer, have a worldwide fame. Hence, Yerevan is home to many leading enterprises of Armenia and the Caucasus for the production of alcoholic beverages, such as the Yerevan Ararat Brandy Factory, Yerevan Brandy Company, Yerevan Champagne Wines Factory, "Beer of Yerevan" (Kilikia Beer) brewery, Armco Brandy Factory, Proshyan Brandy Factory and Astafian Wine-Brandy Factory. The 2 tobacco producers in Yerevan are the "Cigaronne" and "Grand Tabak" companies.
+
+Carpet industry in Armenia has a deeply rooted history with ancient traditions, therefore, carpet production is rather developed in Yerevan with three major factories that also produce hand-made rugs.[146][147][148] The "Megerian Carpet" factory is the leading in this sector.
+
+Other major plants in the city include the "Nairit" chemical and rubber plant, Rusal Armenal aluminum foil mill, "Grand Candy" Armenian-Canadian confectionery manufacturers, "Arcolad" chocolate factory, "Marianna" factory for dairy products, "Talgrig Group" for wheat and flour products, "Shant" ice cream factory, "Crown Chemicals" for paints, "ATMC" travertine mining company, Yerevan Watch Factory "AWI watches", Yerevan Jewellry Plant, and the mineral water factories of "Arzni", "Sil", and "Dilijan Frolova".
+
+Food products include processed meat, all types of canneries, wheat and flour, sweets and chocolate, dried fruits, soft drinks and beverages. Building materials mainly include travertine, crushed stones, asphalt and asphalt concrete.
+
+As an attractive outsourcing location for Western European, Russian and American multinationals, Yerevan headquarters many international companies. It is Armenia's financial hub, being home to the Central Bank of Armenia, the Armenian Stock Exchange (NASDAQ OMX Armenia), as well as the majority of the country's largest commercial banks.[149] As of 2013[update], the city dominates over 85% of the annual total services in Armenia, as well as over 84% of the annual total retail trade.
+
+Many subsidiaries of Russian service companies and banks operate in Yerevan, including Gazprom, Ingo Armenia, Rosgosstrakh and VTB Bank. The ACBA-Credit Agricole is a subsidiary of the French Crédit Agricole, while the HSBC Bank Armenia is also operating in Yerevan.
+
+The construction sector has experienced a significant growth during the 1st decade of the 21st century.[150] Starting from 2000, Yerevan has witnessed a massive construction boom, funded mostly by Armenian millionaires from Russia and the United States, with an extensive and controversial redevelopment process in which many 18th and 19th-century buildings have been demolished and replaced with new buildings. This growth was coupled with a significant increase in real estate prices.[151]
+
+Many major construction projects has been conducted in Yerevan, such as the Northern Avenue and the rehabilitation of Old Yerevan on Aram Street. The Northern Avenue is completed and was opened in 2007, while the Old Yerevan project is still under development. In the past few years, the city centre has also witnessed major road reconstruction, as well as the renovation of the Republic square, funded by the American-Armenian billionaire Kirk Kerkorian. On the other hand, the Argentina-based Armenian businessman Eduardo Eurnekian took over the airport, while the cascade development project was funded by the US based Armenian millionaire Gerard L. Cafesjian.
+
+However, the sector has significantly dropped by the end of the 1st decade of the 21st century, as a result of the global real estate crisis in 2007–09. In 2013, Yerevan dominated over 58% of the annual total construction sector of Armenia.
+
+In February 2017, the urban development committee of the government revealed its plans for the upcoming major construction projects in the city. With a total cost of US$300 million, a new business district will rise at the centre of the city, to replace the current Firdowsi shopping area.[152] The committee has also announced the construction of Noy (Noah) ethnographic residential district at the western vicinity of Kentron District, with an approximate cost of US$100 million.[153]
+
+The location of the city on the shores of Hrazdan river has enabled the production of hydroelectricity. As part of the Sevan–Hrazdan Cascade, three hydroelectric power plants are established within the administrative territory of Yerevan: Kanaker HPP,[154] Yerevan-1 HPP,[155] and Yerevan-3 HPP.[156] The entire plant was privatized in 2003, and is currently owned by RusHydro.[157][158]
+
+The city is also home to the Yerevan Thermal Power Plant, a unique facility in the region for its quality and high technology, situated in the southern part of the city. Originally opened in 1961, a modern plant was built in 2007, furnished with a new gas-steam combined cycled turbine, to generate electric power.[159][160] In March 2017, the construction of a new thermal power plant was launched with an initial investment of US$258 million and an envisaged capacity of 250 megawatts. The power station will be in service in 2019.[161]
+
+As of 2017, Armenia has three mobile phone service providers:
+
+In addition to the mobile network providers, many other small and middle-size companies are also involved in internet services. Access to the Internet in Armenia is largely unfettered. However, according to Article 11 of the Law of the Republic of Armenia on Police, law enforcement has the right to block content to prevent criminal activity.[165]
+
+HayPost is the official national postal operator of Armenia. Based in Yerevan, it currently operates through 900 postal offices across Armenia.[166]
+
+Tourism in Armenia is developing year by year and the capital city of Yerevan is one of the major tourist destinations.[167] The city has a majority of luxury hotels, modern restaurants, bars, pubs and nightclubs. Zvartnots airport has also conducted renovation projects with the growing number of tourists visiting the country. Numerous places in Yerevan are attractive for tourists, such as the dancing fountains of the Republic Square, the State Opera House, the Cascade complex, the ruins of the Urartian city of Erebuni (Arin Berd), the historical site of Karmir Blur (Teishebaini), etc. The largest hotel of the city is the Ani Plaza Hotel. The Armenia Marriott Hotel is located at the Republic Square at the centre of Yerevan, while the Radisson Blu Hotel is located near the Victory Park. Other major chains operating in central Yerevan include the Grand Hotel Yerevan of the Small Luxury Hotels of the World,[168] the Best Western Congress Hotel, the DoubleTree by Hilton, the Hyatt Place, the Ibis Yerevan Center, and The Alexander, a Luxury Collection Hotel of Marriott International.[169]
+
+The location of Yerevan itself, is an inspiring factor for the foreigners to visit the city in order to enjoy the view of the biblical mount of Ararat, as the city lies on the feet of the mountain forming the shape of a Roman amphitheatre.
+
+There are many historical sites, churches and citadels in areas and regions surrounding the city of Yerevan, such as Garni Temple, Zvartnots Cathedral, the monasteries of Khor Virap and Geghard, etc.
+
+Being among the top 10 safest cities in the world, Yerevan has an extensive nightlife scene with a variety of nightclubs,[170] live venues, pedestrian zones, street cafés, jazz cafés, tea houses, casinos, pubs, karaoke clubs and restaurants. Casino Shangri La and Pharaon Complex are among the largest leisure and entertainment centres of the city.
+
+Many world-famous music stars, Russian music celebrities, as well as Armenian singers from diaspora, occasionally perform in concerts in Yerevan.
+
+The Yerevan Zoo founded in 1940, the Yerevan Circus opened in 1956, and the Yerevan Water World opened in 2001, are among the popular entertaining centres in the city.
+
+The Northern Avenue that connects the Opera House with Abovyan street is a popular pedestrian zone in Yerevan with modern residential buildings, business centres, restaurants, bars and cafés. Another popular landmarks is the Yerevan Cascade and the "Cafesjian Sculpture Garden" on Tamanyan Street with its pedestrian zone, featuring many coffee shops, bars, restaurants, and pubs at the sidewalks. The "Cafesjian Center for the Arts" regularly organizes art events throughout the year, including classical music series, traditional folk dance events, and live concerts of jazz, pop and rock music.[171]
+
+As of 2017, Yerevan has three shopping malls: Dalma Garden Mall opened in October 2012, followed by Yerevan Mall in February 2014, and Rossia Mall in March 2016.
+
+International study conducted by Mercer and published in 2019 identified Yerevan to offer higher quality of living, than other capital cities of Transcaucasia.[172][173]
+
+Yerevan is a major educational centre in the region. As of 2017[update], the city is home to more than 250 schools, of which about 210 are state-owned, with 3/4 of them run by the municipality and the rest run by the ministry of education. The rest of the schools (about 40) are privately owned. The municipality also runs 160 kindergartens throughout the city.[174]
+
+The QSI International School, École Française Internationale en Arménie, Ayb School, Mkhitar Sebastatsi Educational Complex and Khoren and Shooshanig Avedisian School are among the prominent international or private schools in Yerevan.
+
+As of 2018[update], around 60 higher education institutions are accredited and licensed to operate in the Republic of Armenia. Yerevan is home to about 50 universities, nearly half of which are public. Yerevan State University, American University of Armenia, Russian-Armenian (Slavonic) University, Yerevan State Medical University and Armenian State Pedagogical University are the top rated universities of Armenia and among the top rated in the region.[175]
+
+Under the Soviet rule, Yerevan has turned into a major centre for science and research. The Armenian National Academy of Sciences is the pioneer of scientific research in Armenia. It was founded in 1943 as the Armenian Branch of the Soviet Academy of Sciences to become the primary body that conducts research and coordinates activities in the fields of science in Armenia. It has many divisions, including Mathematical and Technical Sciences, Physics and Astrophysics, Natural Sciences, Chemistry and Earth Sciences, Armenology and Social Sciences.[176]
+
+After the independence, many new research centres were opened in the city, such as the CANDLE Synchrotron Research Institute (2010),[177] Tumo Center for Creative Technologies (2011),[178] and Nerses Mets Medical Research and Education Center (2013).[179]
+
+Football is the most played and popular sport in Yerevan and the entire country. Yerevan city is home to about a dozen of football clubs competing in the Armenian Premier League and the Armenian First League, with the most successful clubs being Pyunik, Alashkert, Ararat Yerevan, Ararat-Armenia, Urartu and Yerevan.[180]
+
+Hrazdan Stadium in Yerevan is the largest sports venue of Armenia. The 2nd-largest stadium in the city is the Vazgen Sargsyan Republican Stadium which currently serves as the primary home ground of the Armenia national football team.[181]
+
+The Football Academy of Yerevan operated by the Football Federation of Armenia is an up-to-date training academy complex, opened in 2010.[182]
+
+As of 2017, there are around 130 mini-football pitches among the courtyards of the Yerevan neighborhoods, built by the municipal authorities.[183]
+
+Armenia has always excelled in chess with its players being very often among the highest ranked and decorated. The headquarters of the Chess Federation of Armenia is located in the Tigran Petrosian Chess House of Yerevan.[184] The city is home to a large number of chess teams and training schools. In 1996, despite the severe economic conditions in the country, Yerevan hosted the 32nd Chess Olympiad.[185] In 2006, the four members from Yerevan of the Armenian chess team won the 37th Chess Olympiad in Turin and repeated the feat at the 38th Chess Olympiad in Dresden. Armenian won the chess Olympiad for the 3rd time in 2012 in Istanbul. The Yerevan-born leader of the chess national team; Levon Aronian, is one of the top chess players in the world.
+
+Despite the popularity of basketball in Armenia, the country's national team only recently made headlines internationally through winning the 2016 FIBA European Championship for Small Countries. However, the country's best players are diaspora Armenians, mainly from the United States and Russia.
+
+The first ever season of the professional domestic basketball competition of Armenia, known as Armenia Basketball League A, was launched in October 2017 with 7 participating teams. Yerevan is represented by 4 clubs: Engineer Yerevan, FIMA Basketball, BC Grand Sport and BC Urartu.[186]
+
+Tennis is also among the popular sports in Yerevan. Several tennis clubs operate in the city, with many of them founded during the Soviet days. Incourt Tennis Club -founded in 1974– is the largest in the city, with many indoor and outdoor courts.[187] Ararat Tennis Club founded in 1990, is also among the prominent clubs in the city.[188] Tennis clubs are also found within the Yerevan State Sports College of Olympic Reserve since 1971, and the Yerevan Football Academy since 2010.
+
+Sargis Sargsian and Ani Amiraghyan are the most successful tennis players of Armenia.
+
+Armenia has produced many Olympic champions in artistic gymnastics during the Soviet days, such as Hrant Shahinyan, Albert Azaryan and Eduard Azaryan. The success of the Armenian gymnasts in the Olympic competitions has greatly contributed in the popularity of the sport. Thus, many prominent competitors represent the country in the European and World championships, including Artur Davtyan and Harutyun Merdinyan.
+
+Yerevan has many state-owned schools of artistic gymnastics, including the Albert Azaryan School opened in 1964 and the Hrant Shahinyan School opened in 1965.
+
+Karen Demirchyan Sports and Concerts Complex[189] is the largest indoor arena in the city and the entire country. It is mostly used for indoor sport events, including ice hockey and figure skating shows. On the other hand, Dinamo and Mika indoor arenas are the regular venues for domestic and regional competitions of basketball, volleyball, handball and futsal.[190]
+
+Armenia Sports Union (Spartak Sports Union between 1935 and 1999) is a sports society mainly involved in individual Olympic sports, including boxing, weightlifting, athletics, wrestling, taekwondo, table tennis, etc.[191]
+
+The "Yerevan State Sports College of Olympic Reserve" is a large sports and educational complex located in the Malatia-Sebastia District of the city. It was founded in 1971, and is home to individual as well as team sport schools, such as wrestling, boxing, weightlifting, judo, athletics, acrobatic gymnastics, artistic gymnastics, swimming, table tennis, cycling, basketball, volleyball and handball.[192]
+
+In September 2015, the new Olympic Training Complex of Yerevan, locally known as Olympavan, was opened in Davtashen District. It is a state of the art sports complex, with training facilities for most Olympic individual and team sports, as well as water sports. It is also home to the anti-doping medical centre and a hotel designated to accommodate more than 300 athletes.[193]
+
+Equestrian sport was introduced to Armenia in 1953. The Hovik Hayrapetyan Equestrian Centre opened in 2001, occupies an area of 85 hectares at the southern Shengavit District of Yerevan. It is the centre of equestrian sport and horse racing in Armenia.[194]
+
+Golf has been introduced to the citizens of Yerevan in 1999, with the foundation of the Ararat Valley Country Club in the Vahakni neighbourhood of Ajapnyak District. It is the first-ever golf course opened in Armenia as well as the Transcaucasian region.[195]
+
+Arena Bowling and Billiards Club is an up-to-date sports and leisure centre opened in 2004 and located on Mashtots Avenue in central Yerevan.[196]
+
+Cycling as a sport is becoming popular among the young generation. The Yerevan Velodrome is an outdoor track cycling venue with international standard, opened in 2011 to replace the old venue of the Soviet days.[197] Edgar Stepanyan of Armenia became champion of the scratch race in the 2015 junior UEC European Track Championships.[198]
+
+In an attempt to promote figure skating and ice hockey in Armenia, the Irina Rodnina Figure Skating Centre was opened in Yerevan, in December 2015.[199]
+
+Futsal is also among the popular sports in Armenia. Many companies as well as universities have their own teams who participate in the Armenian Futsal Premier League. Currently, Futsal Club Leo based in Yerevan, is considered as the most successful team in the Armenian Futsal Premier League.[180]
+
+Recently, MMA has gained massive popularity in Armenia, being promoted by Armfighting Professional Federation based in Yerevan. It was founded in 2005 by Hayk Ghukasyan and currently runs several branches throughout the provinces of Armenia and Artsakh with more than 2,000 athletes.[200]
+
+With the increased interest in healthy lifestyle and fitness, many large and modern training complexes with indoor and outdoor swimming pools have recently been opened in the city such as the Davit Hambardzumyan Swimming and Diving Olympic School, Orange Fitness Premium Club, DDD Sports Complex, Aqua Land Sports Complex, Gold's Gym, Grand Sport Complex, Reebok Sports Club, and Multi Wellness Sport and Health Center.
+
+The city of Yerevan is member of many international organizations: the International Assembly of CIS Countries' Capitals and Big Cities (MAG), the Black Sea Capitals' Association (BSCA), the International Association of Francophone Mayors (AIMF),[201] the Organization of World Heritage Cities (OWHC), the International Association of Large-scale Communities, and the International Urban Community Lighting Association (LUCI).
+
+Yerevan is twinned with:[202]
+
+Yerevan also cooperates with:[203]

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+Eris typically refers to:
+
+Eris may also refer to:

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+
+
+Eris (minor planet designation 136199 Eris) is the most massive[20] and second-largest known dwarf planet in the Solar System. Eris was discovered in January 2005 by a Palomar Observatory-based team led by Mike Brown, and its discovery was verified later that year. In September 2006 it was named after the goddess of strife and discord. Eris is the ninth-most massive object directly orbiting the Sun, and the sixteenth-most massive overall in the Solar System (including moons). It is also the largest object that has not been visited by a spacecraft. Eris has been measured at 2,326 ± 12 kilometers (1,445.3 ± 7.5 mi) in diameter.[10] Its mass is 0.27 percent as much as the Earth's and 27 percent more than dwarf planet Pluto's,[12][21] though Pluto is slightly larger by volume.[22]
+
+Eris is a trans-Neptunian object (TNO) and a member of a high-eccentricity population known as the scattered disk. It has one known moon, Dysnomia. In February 2016, its distance from the Sun was 96.3 astronomical units (1.441×1010 km; 8.95×109 mi),[17] roughly three times that of Pluto. With the exception of some long-period comets, until 2018 VG18 was discovered on December 17, 2018, Eris and Dysnomia were the most distant known natural objects in the Solar System.[17]
+
+Because Eris appeared to be larger than Pluto, NASA initially described it as the Solar System's tenth planet. This, along with the prospect of other objects of similar size being discovered in the future, motivated the International Astronomical Union (IAU) to define the term planet for the first time. Under the IAU definition approved on August 24, 2006, Eris is a "dwarf planet", along with objects such as Pluto, Ceres, Haumea and Makemake,[23] thereby reducing the number of known planets in the Solar System to eight, the same as before Pluto's discovery in 1930. Observations of a stellar occultation by Eris in 2010 showed that its diameter was 2,326 ± 12 kilometers (1,445.3 ± 7.5 mi), very slightly less than Pluto,[24][25] which was measured by New Horizons as 2,376.6 ± 3.6 kilometers (1,476.8 ± 2.2 mi) in July 2015.[26][27]
+
+Eris was discovered by the team of Mike Brown, Chad Trujillo, and David Rabinowitz[2] on January 5, 2005, from images taken on October 21, 2003.[28] The discovery was announced on July 29, 2005, the same day as Makemake and two days after Haumea,[29] due in part to events that would later lead to controversy about Haumea. The search team had been systematically scanning for large outer Solar System bodies for several years, and had been involved in the discovery of several other large TNOs, including 50000 Quaoar, 90482 Orcus, and 90377 Sedna.[30]
+
+Routine observations were taken by the team on October 21, 2003, using the 1.2 m Samuel Oschin Schmidt telescope at Palomar Observatory, California, but the image of Eris was not discovered at that point due to its very slow motion across the sky: The team's automatic image-searching software excluded all objects moving at less than 1.5 arcseconds per hour to reduce the number of false positives returned.[28] When Sedna was discovered in 2003, it was moving at 1.75 arcsec/h, and in light of that the team reanalyzed their old data with a lower limit on the angular motion, sorting through the previously excluded images by eye. In January 2005, the re-analysis revealed Eris's slow motion against the background stars.[28]
+
+Follow-up observations were then carried out to make a preliminary determination of Eris's orbit, which allowed the object's distance to be estimated.[28] The team had planned to delay announcing their discoveries of the bright objects Eris and Makemake until further observations and calculations were complete, but announced them both on July 29 when the discovery of another large TNO they had been tracking, Haumea, was controversially announced on July 27 by a different team in Spain.[2]
+
+Precovery images of Eris have been identified back to September 3, 1954.[3]
+
+More observations released in October 2005 revealed that Eris has a moon, later named Dysnomia. Observations of Dysnomia's orbit permitted scientists to determine the mass of Eris, which in June 2007 they calculated to be (1.66±0.02)×1022 kg,[12] 27%±2% greater than Pluto's.
+
+Eris is named after the Greek goddess Eris (Greek Ἔρις), a personification of strife and discord.[31] The name was proposed by the Caltech Team on September 6, 2006, and it was assigned on September 13, 2006,[32] following an unusually long period in which the object was known by the provisional designation 2003 UB313, which was granted automatically by the IAU under their naming protocols for minor planets.
+
+Like the moons Io and Mimas, and for the same reason, the name Eris has two competing pronunciations, with a 'long' and a 'short' e.[33] The literary English pronunciation of the goddess is /ˈɪərɪs/ with a long e.[6] However, Brown and his students[34] use something closer to Latin and Greek, /ˈɛrɪs/ with a short e (ignoring the Great Vowel Shift that affects Classical names in English).[7]
+
+The Greek and Latin oblique stem of the name is Erid-,[35] as can be seen in Italian Eride and Russian Эрида Erida, so the adjective in English is Eridian /ɛˈrɪdiən/.[8][9]
+
+Due to uncertainty over whether the object would be classified as a planet or a minor planet, because different nomenclature procedures apply to these different classes of objects,[36] the decision on what to name the object had to wait until after the August 24, 2006 IAU ruling.[37] As a result, for a time the object became known to the wider public as Xena.
+
+"Xena" was an informal name used internally by the discovery team. It was inspired by the title character of the television series Xena: Warrior Princess. The discovery team had reportedly saved the nickname "Xena" for the first body they discovered that was larger than Pluto. According to Brown,
+
+We chose it since it started with an X (planet "X"), it sounds mythological (OK, so it's TV mythology, but Pluto is named after a cartoon, right?),[e] and (this part is actually true) we've been working to get more female deities out there (e.g. Sedna). Also, at the time, the TV show was still on TV, which shows you how long we've been searching![39]
+
+"We assumed [that] a real name would come out fairly quickly, [but] the process got stalled", Mike Brown said in an interview:
+
+One reporter [Ken Chang][40] called me up from The New York Times who happened to have been a friend of mine from college, [and] I was a little less guarded with him than I am with the normal press. He asked me, "What's the name you guys proposed?" and I said, "Well, I'm not going to tell." And he said, "Well, what do you guys call it when you're just talking amongst yourselves?" ... As far as I remember this was the only time I told anybody this in the press, and then it got everywhere, which I only sorta felt bad about—I kinda like the name.[41]
+
+According to science writer Govert Schilling, Brown initially wanted to call the object "Lila", after a concept in Hindu mythology that described the cosmos as the outcome of a game played by Brahman.[30] The name was very similar to "Lilah", the name of Brown's newborn daughter. Brown was mindful of not making his name public before it had been officially accepted. He had done so with Sedna a year previously, and had been heavily criticized. However, no objection was raised to the Sedna name other than the breach of protocol, and no competing names were suggested for Sedna.[42]
+
+He listed the address of his personal web page announcing the discovery as /~mbrown/planetlila and in the chaos following the controversy over the discovery of Haumea, forgot to change it. Rather than needlessly anger more of his fellow astronomers, he simply said that the webpage had been named for his daughter and dropped "Lila" from consideration.[30]
+
+Brown had also speculated that Persephone, the wife of the god Pluto, would be a good name for the object.[2] The name had been used several times in science fiction,[43] and was popular with the public, having handily won a poll conducted by New Scientist magazine ("Xena", despite only being a nickname, came fourth).[44] This was not possible once the object was classified as a dwarf planet, because there is already an asteroid with that name, 399 Persephone.[2]
+
+With the dispute resolved, the discovery team proposed Eris on September 6, 2006. On September 13, 2006 this name was accepted as the official name by the IAU.[45][46] Brown decided that, because the object had been considered a planet for so long, it deserved a name from Greek or Roman mythology, like the other planets. The asteroids had taken the vast majority of Graeco-Roman names. Eris, whom Brown described as his favorite goddess, had fortunately escaped inclusion.[41] "Eris caused strife and discord by causing quarrels among people", said Brown in 2006, "and that's what this one has done too".[47]
+
+Eris is a trans-Neptunian dwarf planet (plutoid).[48] Its orbital characteristics more specifically categorize it as a scattered-disk object (SDO), or a TNO that has been "scattered" from the Kuiper belt into more-distant and unusual orbits following gravitational interactions with Neptune as the Solar System was forming. Although its high orbital inclination is unusual among the known SDOs, theoretical models suggest that objects that were originally near the inner edge of the Kuiper belt were scattered into orbits with higher inclinations than objects from the outer belt.[49]
+
+Because Eris was initially thought to be larger than Pluto, it was described as the "tenth planet" by NASA and in media reports of its discovery.[50] In response to the uncertainty over its status, and because of ongoing debate over whether Pluto should be classified as a planet, the IAU delegated a group of astronomers to develop a sufficiently precise definition of the term planet to decide the issue. This was announced as the IAU's Definition of a Planet in the Solar System, adopted on August 24, 2006. At this time, both Eris and Pluto were classified as dwarf planets, a category distinct from the new definition of planet.[51] Brown has since stated his approval of this classification.[52] The IAU subsequently added Eris to its Minor Planet Catalogue, designating it (136199) Eris.[37]
+
+Eris has an orbital period of 559 years.[17] Its maximum possible distance from the Sun (aphelion) is 97.65 AU, and its closest (perihelion) is 37.91 AU.[17] It came to perihelion between 1698[5] and 1699,[53] to aphelion around 1977,[53] and will return to perihelion around 2256[53] to 2258.[54] Unlike the eight planets, whose orbits all lie roughly in the same plane as the Earth's, Eris's orbit is highly inclined: It is tilted at an angle of about 44 degrees to the ecliptic.[3] When discovered, Eris and its moon were the most distant known objects in the Solar System, apart from long-period comets and space probes.[2][55] It retained this distinction until the discovery of 2018 VG18 in 2018.[56]
+
+As of 2008 there were approximately forty known TNOs, most notably 2006 SQ372, 2000 OO67 and Sedna, that are currently closer to the Sun than Eris even though their semimajor axis is larger than that of Eris (67.8 AU).[4]
+
+Eris's orbit is highly eccentric, and brings Eris to within 37.9 AU of the Sun, a typical perihelion for scattered objects.[57]  This is within the orbit of Pluto, but still safe from direct interaction with Neptune (~37 AU).[58] Pluto, on the other hand, like other plutinos, follows a less inclined and less eccentric orbit and, protected by orbital resonance, can cross Neptune's orbit.[59] In about 800 years, Eris will be closer to the Sun than Pluto for some time (see the graph at the left).
+
+As of 2007, Eris has an apparent magnitude of 18.7, making it bright enough to be detectable to some amateur telescopes.[60] A 200-millimeter (7.9 in) telescope with a CCD can detect Eris under favorable conditions.[f] The reason it had not been noticed until now is its steep orbital inclination; searches for large outer Solar System objects tend to concentrate on the ecliptic plane, where most bodies are found.[61]
+
+Because of the high inclination of its orbit, Eris passes through only a few constellations of the traditional Zodiac; it is now in the constellation Cetus. It was in Sculptor from 1876 until 1929 and Phoenix from roughly 1840 until 1875. In 2036 it will enter Pisces and stay there until 2065, when it will enter Aries.[53] It will then move into the northern sky, entering Perseus in 2128 and Camelopardalis (where it will reach its northernmost declination) in 2173.
+
+In November 2010, Eris was the subject of one of the most distant stellar occultations yet from Earth.[11] Preliminary data from this event cast doubt on previous size estimates.[11] The teams announced their final results from the occultation in October 2011, with an estimated diameter of 2326±12 km.[10]
+
+This makes Eris a little smaller than Pluto by area and diameter, which is 2372±4 km across, although Eris is more massive. It also indicates an albedo of 0.96, higher than that of any other large body in the Solar System except Enceladus.[10] It is speculated that the high albedo is due to the surface ices being replenished because of temperature fluctuations as Eris's eccentric orbit takes it closer and farther from the Sun.[19]
+
+The mass of Eris can be calculated with much greater precision. Based on the currently accepted value for Dysnomia's period—15.774 days[12][64]—Eris is 27 percent more massive than Pluto. Using the 2011 occultation results, Eris has a density of 2.52±0.07 g/cm3,[b] substantially denser than Pluto, and thus must be composed largely of rocky materials.[10]
+
+Models of internal heating via radioactive decay suggest that Eris could have an internal ocean of liquid water at the mantle–core boundary.[65]
+
+In July 2015, after nearly ten years of Eris being considered the ninth-largest object known to directly orbit the sun, close-up imagery from the New Horizons mission more accurately determined Pluto's volume to be slightly larger than Eris's, rather than slightly smaller as previously thought.[66] Eris is now the tenth-largest object known to directly orbit the sun by volume, but remains the ninth-largest by mass.
+
+The discovery team followed up their initial identification of Eris with spectroscopic observations made at the 8 m Gemini North Telescope in Hawaii on January 25, 2005. Infrared light from the object revealed the presence of methane ice, indicating that the surface may be similar to that of Pluto, which at the time was the only TNO known to have surface methane, and of Neptune's moon Triton, which also has methane on its surface.[67]
+
+Due to Eris's distant eccentric orbit, its surface temperature is estimated to vary between about 30 and 56 K (−243.2 and −217.2 °C).[2]
+
+Unlike the somewhat reddish Pluto and Triton, Eris appears almost white.[2] Pluto's reddish color is thought to be due to deposits of tholins on its surface, and where these deposits darken the surface, the lower albedo leads to higher temperatures and the evaporation of methane deposits. In contrast, Eris is far enough from the Sun that methane can condense onto its surface even where the albedo is low. The condensation of methane uniformly over the surface reduces any albedo contrasts and would cover up any deposits of red tholins.[28]
+
+Even though Eris can be up to three times farther from the Sun than Pluto, it approaches close enough that some of the ices on the surface might warm enough to sublime. Because methane is highly volatile, its presence shows either that Eris has always resided in the distant reaches of the Solar System, where it is cold enough for methane ice to persist, or that the celestial body has an internal source of methane to replenish gas that escapes from its atmosphere. This contrasts with observations of another discovered TNO, Haumea, which reveal the presence of water ice but not methane.[68]
+
+In 2005, the adaptive optics team at the Keck telescopes in Hawaii carried out observations of the four brightest TNOs (Pluto, Makemake, Haumea, and Eris), using the newly commissioned laser guide star adaptive optics system.[70] Images taken on September 10 revealed a moon in orbit around Eris. In keeping with the "Xena" nickname already in use for Eris, Brown's team nicknamed the moon "Gabrielle", after the television warrior princess' sidekick. When Eris received its official name from the IAU, the moon received the name Dysnomia, after the Greek goddess of lawlessness who was Eris's daughter. Brown says he picked it for similarity to his wife's name, Diane. The name also retains an oblique reference to Eris's old informal name Xena, portrayed on TV by Lucy Lawless.[71]
+
+In the 2010s, on the heels of the successful Pluto flyby there were multiple studies for follow-on missions to explore the Kuiper belt, and Eris was evaluated among the candidates.[72] It was calculated that a flyby mission to Eris could take 24.66 years using a Jupiter gravity assist, based on launch dates of April 3, 2032 or April 7, 2044. Eris would be 92.03 or 90.19 AU from the Sun when the spacecraft arrives.[73]
+
+Solar System → Local Interstellar Cloud → Local Bubble → Gould Belt → Orion Arm → Milky Way → Milky Way subgroup → Local Group → Local Sheet → Virgo Supercluster → Laniakea Supercluster → Observable universe → UniverseEach arrow (→) may be read as "within" or "part of".
+

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+Ernesto "Che" Guevara (/tʃeɪ ɡəˈvɑːrə/;[3] Spanish: [ˈtʃe ɣeˈβaɾa];[4] 14 June 1928[5] – 9 October 1967) was an Argentine Marxist revolutionary, physician, author, guerrilla leader, diplomat, and military theorist. A major figure of the Cuban Revolution, his stylized visage has become a ubiquitous countercultural symbol of rebellion and global insignia in popular culture.[6]
+
+As a young medical student, Guevara traveled throughout South America and was radicalized by the poverty, hunger, and disease he witnessed.[7][8] His burgeoning desire to help overturn what he saw as the capitalist exploitation of Latin America by the United States prompted his involvement in Guatemala's social reforms under President Jacobo Árbenz, whose eventual CIA-assisted overthrow at the behest of the United Fruit Company solidified Guevara's political ideology.[7] Later in Mexico City, Guevara met Raúl and Fidel Castro, joined their 26th of July Movement, and sailed to Cuba aboard the yacht Granma with the intention of overthrowing U.S.-backed Cuban dictator Fulgencio Batista.[9] Guevara soon rose to prominence among the insurgents, was promoted to second in command and played a pivotal role in the victorious two-year guerrilla campaign that deposed the Batista regime.[10]
+
+Following the Cuban Revolution, Guevara performed a number of key roles in the new government. These included reviewing the appeals and firing squads for those convicted as war criminals during the revolutionary tribunals,[11] instituting agrarian land reform as minister of industries, helping spearhead a successful nationwide literacy campaign, serving as both national bank president and instructional director for Cuba's armed forces, and traversing the globe as a diplomat on behalf of Cuban socialism. Such positions also allowed him to play a central role in training the militia forces who repelled the Bay of Pigs Invasion,[12] and bringing Soviet nuclear-armed ballistic missiles to Cuba, which preceded the 1962 Cuban Missile Crisis.[13] Additionally, Guevara was a prolific writer and diarist, composing a seminal guerrilla warfare manual, along with a best-selling memoir about his youthful continental motorcycle journey. His experiences and studying of Marxism–Leninism led him to posit that the Third World's underdevelopment and dependence was an intrinsic result of imperialism, neocolonialism and monopoly capitalism, with the only remedy being proletarian internationalism and world revolution.[14][15] Guevara left Cuba in 1965 to foment revolution abroad, first unsuccessfully in Congo-Kinshasa and later in Bolivia, where he was captured by CIA-assisted Bolivian forces and summarily executed.[16]
+
+Guevara remains both a revered and reviled historical figure, polarized in the collective imagination in a multitude of biographies, memoirs, essays, documentaries, songs, and films. As a result of his perceived martyrdom, poetic invocations for class struggle, and desire to create the consciousness of a "new man" driven by moral rather than material incentives,[17] Guevara has evolved into a quintessential icon of various leftist movements. In contrast, his ideological critics on the right accuse him of authoritarianism and sanctifying violence against his political opponents. Despite disagreements on his legacy, Time magazine named him one of the 100 most influential people of the 20th century,[18] while an Alberto Korda photograph of him, titled Guerrillero Heroico, was cited by the Maryland Institute College of Art as "the most famous photograph in the world".[19]
+
+Ernesto Guevara was born to Ernesto Guevara Lynch and Celia de la Serna y Llosa, on 14 June 1928,[5] in Rosario, Argentina. He was the eldest of five children in a middle-class Argentine family of Spanish (including Basque and Cantabrian) descent, as well as Irish by means of his patrilineal ancestor Patrick Lynch.[20][21][22] Although Guevara's legal name on his birth certificate was "Ernesto Guevara", his name sometimes appears with "de la Serna" and/or "Lynch" accompanying it.[23] Referring to Che's "restless" nature, his father declared "the first thing to note is that in my son's veins flowed the blood of the Irish rebels".[24]
+
+Very early on in life, Ernestito (as he was then called) developed an "affinity for the poor".[25] Growing up in a family with leftist leanings, Guevara was introduced to a wide spectrum of political perspectives even as a boy.[26] His father, a staunch supporter of Republicans from the Spanish Civil War, often hosted many veterans from the conflict in the Guevara home.[27]
+
+Despite suffering crippling bouts of acute asthma that were to afflict him throughout his life, he excelled as an athlete, enjoying swimming, football, golf, and shooting, while also becoming an "untiring" cyclist.[28][29] He was an avid rugby union player,[30] and played at fly-half for Club Universitario de Buenos Aires.[31] His rugby playing earned him the nickname "Fuser"—a contraction of El Furibundo (raging) and his mother's surname, de la Serna—for his aggressive style of play.[32]
+
+Guevara learned chess from his father, and began participating in local tournaments by the age of 12. During adolescence and throughout his life he was passionate about poetry, especially that of Pablo Neruda, John Keats, Antonio Machado, Federico García Lorca, Gabriela Mistral, César Vallejo, and Walt Whitman.[33] He could also recite Rudyard Kipling's "If—" and José Hernández's Martín Fierro by heart.[33] The Guevara home contained more than 3,000 books, which allowed Guevara to be an enthusiastic and eclectic reader, with interests including Karl Marx, William Faulkner, André Gide, Emilio Salgari and Jules Verne.[34] Additionally, he enjoyed the works of Jawaharlal Nehru, Franz Kafka, Albert Camus, Vladimir Lenin and Jean-Paul Sartre; as well as Anatole France, Friedrich Engels, H. G. Wells and Robert Frost.[35]
+
+As he grew older, he developed an interest in the Latin American writers Horacio Quiroga, Ciro Alegría, Jorge Icaza, Rubén Darío and Miguel Asturias.[35] Many of these authors' ideas he cataloged in his own handwritten notebooks of concepts, definitions, and philosophies of influential intellectuals. These included composing analytical sketches of Buddha and Aristotle, along with examining Bertrand Russell on love and patriotism, Jack London on society and Nietzsche on the idea of death. Sigmund Freud's ideas fascinated him as he quoted him on a variety of topics from dreams and libido to narcissism and the Oedipus complex.[35] His favorite subjects in school included philosophy, mathematics, engineering, political science, sociology, history and archaeology.[36][37]
+
+Years later, a declassified CIA 'biographical and personality report' dated 13 February 1958 made note of Guevara's wide range of academic interests and intellect, describing him as "quite well read" while adding that "Che is fairly intellectual for a Latino."[38]
+
+In 1948, Guevara entered the University of Buenos Aires to study medicine. His "hunger to explore the world"[40] led him to intersperse his collegiate pursuits with two long introspective journeys that fundamentally changed the way he viewed himself and the contemporary economic conditions in Latin America. The first expedition in 1950 was a 4,500-kilometer (2,800 mi) solo trip through the rural provinces of northern Argentina on a bicycle on which he installed a small engine.[41] This was followed in 1951 by a nine-month, 8,000-kilometer (5,000 mi) continental motorcycle trek through part of South America. For the latter, he took a year off from his studies to embark with his friend Alberto Granado, with the final goal of spending a few weeks volunteering at the San Pablo leper colony in Peru, on the banks of the Amazon River.[42]
+
+In Chile, Guevara found himself enraged by the working conditions of the miners in Anaconda's Chuquicamata copper mine and moved by his overnight encounter in the Atacama Desert with a persecuted communist couple who did not even own a blanket, describing them as "the shivering flesh-and-blood victims of capitalist exploitation".[43] Additionally, on the way to Machu Picchu high in the Andes, he was struck by the crushing poverty of the remote rural areas, where peasant farmers worked small plots of land owned by wealthy landlords.[44] Later on his journey, Guevara was especially impressed by the camaraderie among those living in a leper colony, stating, "The highest forms of human solidarity and loyalty arise among such lonely and desperate people."[44] Guevara used notes taken during this trip to write an account, titled The Motorcycle Diaries, which later became a New York Times best-seller,[45] and was adapted into a 2004 award-winning film of the same name.
+
+—George Galloway, British politician[46]
+
+The journey took Guevara through Argentina, Chile, Peru, Ecuador, Colombia, Venezuela, Panama, and Miami, Florida, for 20 days,[47] before returning home to Buenos Aires. By the end of the trip, he came to view Latin America not as a collection of separate nations, but as a single entity requiring a continent-wide liberation strategy. His conception of a borderless, united Hispanic America sharing a common Latino heritage was a theme that recurred prominently during his later revolutionary activities. Upon returning to Argentina, he completed his studies and received his medical degree in June 1953, making him officially "Dr. Ernesto Guevara".[48][49]
+
+Guevara later remarked that through his travels in Latin America, he came in "close contact with poverty, hunger and disease" along with the "inability to treat a child because of lack of money" and "stupefaction provoked by the continual hunger and punishment" that leads a father to "accept the loss of a son as an unimportant accident". Guevara cited these experiences as convincing him that in order to "help these people", he needed to leave the realm of medicine and consider the political arena of armed struggle.[7]
+
+On 7 July 1953, Guevara set out again, this time to Bolivia, Peru, Ecuador, Panama, Costa Rica, Nicaragua, Honduras and El Salvador. On 10 December 1953, before leaving for Guatemala, Guevara sent an update to his Aunt Beatriz from San José, Costa Rica. In the letter Guevara speaks of traversing the dominion of the United Fruit Company, a journey which convinced him that the Company's capitalist system was a terrible one.[50] This affirmed indignation carried the more aggressive tone he adopted in order to frighten his more Conservative relatives, and ends with Guevara swearing on an image of the then recently deceased Joseph Stalin, not to rest until these "octopuses have been vanquished".[51] Later that month, Guevara arrived in Guatemala where President Jacobo Árbenz Guzmán headed a democratically elected government that, through land reform and other initiatives, was attempting to end the latifundia system. To accomplish this, President Árbenz had enacted a major land reform program, where all uncultivated portions of large land holdings were to be expropriated and redistributed to landless peasants. The biggest land owner, and one most affected by the reforms, was the United Fruit Company, from which the Árbenz government had already taken more than 225,000 acres (91,000 ha) of uncultivated land.[52] Pleased with the road the nation was heading down, Guevara decided to settle down in Guatemala so as to "perfect himself and accomplish whatever may be necessary in order to become a true revolutionary."[53]
+
+In Guatemala City, Guevara sought out Hilda Gadea Acosta, a Peruvian economist who was well-connected politically as a member of the left-leaning Alianza Popular Revolucionaria Americana (APRA, American Popular Revolutionary Alliance). She introduced Guevara to a number of high-level officials in the Arbenz government. Guevara then established contact with a group of Cuban exiles linked to Fidel Castro through the 26 July 1953 attack on the Moncada Barracks in Santiago de Cuba. During this period, he acquired his famous nickname, due to his frequent use of the Argentine filler syllable che (a multi-purpose discourse marker, like the syllable "eh" in Canadian English).[54] During his time in Guatemala, Guevara was helped by other Central American exiles, one of whom, Helena Leiva de Holst, provided him with food and lodging,[55] discussed her travels to study Marxism in Russia and China,[56] and to whom, Guevara dedicated a poem, "Invitación al camino".[57]
+
+In May 1954, a shipment of infantry and light artillery weapons was dispatched from Communist Czechoslovakia for the Arbenz Government and arrived in Puerto Barrios.[58] As a result, the United States government—which since 1953 had been tasked by President Eisenhower to remove Arbenz from power in the multifaceted CIA operation code-named PBSUCCESS—responded by saturating Guatemala with anti-Arbenz propaganda through radio and dropped leaflets, and began bombing raids using unmarked airplanes.[59] The United States also sponsored a force of several hundred Guatemalan refugees and mercenaries who were headed by Castillo Armas to help remove the Arbenz government. On 27 June, Arbenz decided to resign.[60]  This allowed Armas and his CIA-assisted forces to march into Guatemala City and establish a military junta, which elected Armas as President on 7 July.[61] Consequently, the Armas regime then consolidated power by rounding up and executing suspected communists,[62] while crushing the previously flourishing labor unions[63] and reversing the previous agrarian reforms.[64]
+
+Guevara himself was eager to fight on behalf of Arbenz and joined an armed militia organized by the Communist Youth for that purpose, but frustrated with the group's inaction, he soon returned to medical duties. Following the coup, he again volunteered to fight, but soon after, Arbenz took refuge in the Mexican Embassy and told his foreign supporters to leave the country. Guevara's repeated calls to resist were noted by supporters of the coup, and he was marked for murder.[65] After Hilda Gadea was arrested, Guevara sought protection inside the Argentine consulate, where he remained until he received a safe-conduct pass some weeks later and made his way to Mexico.[66]
+
+The overthrow of the Arbenz regime and establishment of the right-wing Armas dictatorship cemented Guevara's view of the United States as an imperialist power that opposed and attempted to destroy any government that sought to redress the socioeconomic inequality endemic to Latin America and other developing countries.[53] In speaking about the coup, Guevara stated:
+
+The last Latin American revolutionary democracy – that of Jacobo Arbenz – failed as a result of the cold premeditated aggression carried out by the United States. Its visible head was the Secretary of State John Foster Dulles, a man who, through a rare coincidence, was also a stockholder and attorney for the United Fruit Company.[65]
+
+Guevara's conviction that Marxism achieved through armed struggle and defended by an armed populace was the only way to rectify such conditions was thus strengthened.[67] Gadea wrote later, "It was Guatemala which finally convinced him of the necessity for armed struggle and for taking the initiative against imperialism. By the time he left, he was sure of this."[68]
+
+Guevara arrived in Mexico City on 21 September 1954, and worked in the allergy section of the General Hospital and at the Hospital Infantil de Mexico.[69][70] In addition he gave lectures on medicine at the Faculty of Medicine in the National Autonomous University of Mexico and worked as a news photographer for Latina News Agency.[71][72] His first wife Hilda notes in her memoir My Life with Che, that for a while, Guevara considered going to work as a doctor in Africa and that he continued to be deeply troubled by the poverty around him.[73] In one instance, Hilda describes Guevara's obsession with an elderly washerwoman whom he was treating, remarking that he saw her as "representative of the most forgotten and exploited class". Hilda later found a poem that Che had dedicated to the old woman, containing "a promise to fight for a better world, for a better life for all the poor and exploited".[73]
+
+During this time he renewed his friendship with Ñico López and the other Cuban exiles whom he had met in Guatemala. In June 1955, López introduced him to Raúl Castro, who subsequently introduced him to his older brother, Fidel Castro, the revolutionary leader who had formed the 26th of July Movement and was now plotting to overthrow the dictatorship of Fulgencio Batista. During a long conversation with Fidel on the night of their first meeting, Guevara concluded that the Cuban's cause was the one for which he had been searching and before daybreak he had signed up as a member of the July 26 Movement.[74] Despite their "contrasting personalities", from this point on Che and Fidel began to foster what dual biographer Simon Reid-Henry deemed a "revolutionary friendship that would change the world", as a result of their coinciding commitment to anti-imperialism.[75]
+
+By this point in Guevara's life, he deemed that U.S.-controlled conglomerates installed and supported repressive regimes around the world. In this vein, he considered Batista a "U.S. puppet whose strings needed cutting".[76] Although he planned to be the group's combat medic, Guevara participated in the military training with the members of the Movement. The key portion of training involved learning hit and run tactics of guerrilla warfare. Guevara and the others underwent arduous 15-hour marches over mountains, across rivers, and through the dense undergrowth, learning and perfecting the procedures of ambush and quick retreat. From the start Guevara was Alberto Bayo's "prize student" among those in training, scoring the highest on all of the tests given.[77] At the end of the course, he was called "the best guerrilla of them all" by their instructor, General Bayo.[78]
+
+Guevara then married Gadea in Mexico in September 1955, before embarking on his plan to assist in the liberation of Cuba.[79]
+
+The first step in Castro's revolutionary plan was an assault on Cuba from Mexico via the Granma, an old, leaky cabin cruiser. They set out for Cuba on 25 November 1956. Attacked by Batista's military soon after landing, many of the 82 men were either killed in the attack or executed upon capture; only 22 found each other afterwards.[80] During this initial bloody confrontation Guevara laid down his medical supplies and picked up a box of ammunition dropped by a fleeing comrade, proving to be a symbolic moment in Che's life.[81]
+
+Only a small band of revolutionaries survived to re-group as a bedraggled fighting force deep in the Sierra Maestra mountains, where they received support from the urban guerrilla network of Frank País, 26 July Movement, and local campesinos. With the group withdrawn to the Sierra, the world wondered whether Castro was alive or dead until early 1957 when the interview by Herbert Matthews appeared in The New York Times. The article presented a lasting, almost mythical image for Castro and the guerrillas. Guevara was not present for the interview, but in the coming months he began to realize the importance of the media in their struggle. Meanwhile, as supplies and morale diminished, and with an allergy to mosquito bites which resulted in agonizing walnut-sized cysts on his body,[82] Guevara considered these "the most painful days of the war".[83]
+
+During Guevara's time living hidden among the poor subsistence farmers of the Sierra Maestra mountains, he discovered that there were no schools, no electricity, minimal access to healthcare, and more than 40 percent of the adults were illiterate.[84] As the war continued, Guevara became an integral part of the rebel army and "convinced Castro with competence, diplomacy and patience".[10] Guevara set up factories to make grenades, built ovens to bake bread, and organized schools to teach illiterate campesinos to read and write.[10] Moreover, Guevara established health clinics, workshops to teach military tactics, and a newspaper to disseminate information.[85] The man whom Time dubbed three years later "Castro's brain" at this point was promoted by Fidel Castro to Comandante (commander) of a second army column.[10]
+
+As second in command, Guevara was a harsh disciplinarian who sometimes shot defectors. Deserters were punished as traitors, and Guevara was known to send squads to track those seeking to go AWOL.[86] As a result, Guevara became feared for his brutality and ruthlessness.[87] During the guerrilla campaign, Guevara was also responsible for the summary executions of a number of men accused of being informers, deserters or spies.[88] In his diaries, Guevara described the first such execution of Eutímio Guerra, a peasant army guide who admitted treason when it was discovered he accepted the promise of ten thousand pesos for repeatedly giving away the rebel's position for attack by the Cuban air force.[89] Such information also allowed Batista's army to burn the homes of peasants sympathetic to the revolution.[89] Upon Guerra's request that they "end his life quickly",[89] Che stepped forward and shot him in the head, writing "The situation was uncomfortable for the people and for Eutimio so I ended the problem giving him a shot with a .32 pistol in the right side of the brain, with exit orifice in the right temporal [lobe]."[90] His scientific notations and matter-of-fact description, suggested to one biographer a "remarkable detachment to violence" by that point in the war.[90] Later, Guevara published a literary account of the incident, titled "Death of a Traitor", where he transfigured Eutimio's betrayal and pre-execution request that the revolution "take care of his children", into a "revolutionary parable about redemption through sacrifice".[90]
+
+Although he maintained a demanding and harsh disposition, Guevara also viewed his role of commander as one of a teacher, entertaining his men during breaks between engagements with readings from the likes of Robert Louis Stevenson, Miguel de Cervantes, and Spanish lyric poets.[91] Together with this role, and inspired by José Martí's principle of "literacy without borders", Guevara further ensured that his rebel fighters made daily time to teach the uneducated campesinos with whom they lived and fought to read and write, in what Guevara termed the "battle against ignorance".[84] Tomás Alba, who fought under Guevara's command, later stated that "Che was loved, in spite of being stern and demanding. We would (have) given our life for him."[92]
+
+His commanding officer Fidel Castro described Guevara as intelligent, daring, and an exemplary leader who "had great moral authority over his troops".[93] Castro further remarked that Guevara took too many risks, even having a "tendency toward foolhardiness".[94] Guevara's teenage lieutenant, Joel Iglesias, recounts such actions in his diary, noting that Guevara's behavior in combat even brought admiration from the enemy. On one occasion Iglesias recounts the time he had been wounded in battle, stating "Che ran out to me, defying the bullets, threw me over his shoulder, and got me out of there. The guards didn't dare fire at him ... later they told me he made a great impression on them when they saw him run out with his pistol stuck in his belt, ignoring the danger, they didn't dare shoot."[95]
+
+Guevara was instrumental in creating the clandestine radio station Radio Rebelde (Rebel Radio) in February 1958, which broadcast news to the Cuban people with statements by 26 July movement, and provided radiotelephone communication between the growing number of rebel columns across the island. Guevara had apparently been inspired to create the station by observing the effectiveness of CIA supplied radio in Guatemala in ousting the government of Jacobo Árbenz Guzmán.[96]
+
+To quell the rebellion, Cuban government troops began executing rebel prisoners on the spot, and regularly rounded up, tortured, and shot civilians as a tactic of intimidation.[97] By March 1958, the continued atrocities carried out by Batista's forces led the United States to stop selling arms to the Cuban government.[85] Then in late July 1958, Guevara played a critical role in the Battle of Las Mercedes by using his column to halt a force of 1,500 men called up by Batista's General Cantillo in a plan to encircle and destroy Castro's forces. Years later, Major Larry Bockman of the United States Marine Corps analyzed and described Che's tactical appreciation of this battle as "brilliant".[98] During this time Guevara also became an "expert" at leading hit-and-run tactics against Batista's army, and then fading back into the countryside before the army could counterattack.[99]
+
+As the war extended, Guevara led a new column of fighters dispatched westward for the final push towards Havana. Travelling by foot, Guevara embarked on a difficult 7-week march, only travelling at night to avoid an ambush and often not eating for several days.[100] In the closing days of December 1958, Guevara's task was to cut the island in half by taking Las Villas province. In a matter of days he executed a series of "brilliant tactical victories" that gave him control of all but the province's capital city of Santa Clara.[100] Guevara then directed his "suicide squad" in the attack on Santa Clara, which became the final decisive military victory of the revolution.[101][102] In the six weeks leading up to the battle, there were times when his men were completely surrounded, outgunned, and overrun. Che's eventual victory despite being outnumbered 10:1 remains in the view of some observers a "remarkable tour de force in modern warfare".[103]
+
+Radio Rebelde broadcast the first reports that Guevara's column had taken Santa Clara on New Year's Eve 1958. This contradicted reports by the heavily controlled national news media, which had at one stage reported Guevara's death during the fighting. At 3 am on 1 January 1959, upon learning that his generals were negotiating a separate peace with Guevara, Fulgencio Batista boarded a plane in Havana and fled for the Dominican Republic, along with an amassed "fortune of more than $300,000,000 through graft and payoffs".[104] The following day on 2 January, Guevara entered Havana to take final control of the capital.[105] Fidel Castro took six more days to arrive, as he stopped to rally support in several large cities on his way to rolling victoriously into Havana on 8 January 1959. The final death toll from the two years of revolutionary fighting was 2,000 people.[106]
+
+In mid-January 1959, Guevara went to live at a summer villa in Tarará to recover from a violent asthma attack.[107] While there he started the Tarara Group, a group that debated and formed the new plans for Cuba's social, political, and economic development.[108] In addition, Che began to write his book Guerrilla Warfare while resting at Tarara.[108] In February, the revolutionary government proclaimed Guevara "a Cuban citizen by birth" in recognition of his role in the triumph.[109] When Hilda Gadea arrived in Cuba in late January, Guevara told her that he was involved with another woman, and the two agreed on a divorce,[110] which was finalized on 22 May.[111] On 2 June 1959, he married Aleida March, a Cuban-born member of 26 July movement with whom he had been living since late 1958. Guevara returned to the seaside village of Tarara in June for his honeymoon with Aleida.[112] In total, Guevara had five children from his two marriages.[113]
+
+The first major political crisis arose over what to do with the captured Batista officials who had perpetrated the worst of the repression.[114] During the rebellion against Batista's dictatorship, the general command of the rebel army, led by Fidel Castro, introduced into the territories under its control the 19th-century penal law commonly known as the Ley de la Sierra (Law of the Sierra).[115] This law included the death penalty for serious crimes, whether perpetrated by the Batista regime or by supporters of the revolution. In 1959 the revolutionary government extended its application to the whole of the republic and to those it considered war criminals, captured and tried after the revolution. According to the Cuban Ministry of Justice, this latter extension was supported by the majority of the population, and followed the same procedure as those in the Nuremberg trials held by the Allies after World War II.[116]
+
+To implement a portion of this plan, Castro named Guevara commander of the La Cabaña Fortress prison, for a five-month tenure (2 January through 12 June 1959).[117] Guevara was charged by the new government with purging the Batista army and consolidating victory by exacting "revolutionary justice" against those regarded as traitors, chivatos (informants) or war criminals.[118] As commander of La Cabaña, Guevara reviewed the appeals of those convicted during the revolutionary tribunal process.[11] The tribunals were conducted by 2–3 army officers, an assessor, and a respected local citizen.[119] On some occasions the penalty delivered by the tribunal was death by firing-squad.[120] Raúl Gómez Treto, senior legal advisor to the Cuban Ministry of Justice, has argued that the death penalty was justified in order to prevent citizens themselves from taking justice into their own hands, as had happened twenty years earlier in the anti-Machado rebellion.[121] Biographers note that in January 1959 the Cuban public was in a "lynching mood",[122] and point to a survey at the time showing 93% public approval for the tribunal process.[11] Moreover, a 22 January 1959, Universal Newsreel broadcast in the United States and narrated by Ed Herlihy featured Fidel Castro asking an estimated one million Cubans whether they approved of the executions, and being met with a roaring "¡Si!" (yes).[123] With as many as 20,000 Cubans estimated to have been killed at the hands of Batista's collaborators,[124][125][126][127] and many of the accused war criminals sentenced to death accused of torture and physical atrocities,[11] the newly-empowered government carried out executions, punctuated by cries from the crowds of "¡al paredón!" ([to the] wall!),[114] which biographer Jorge Castañeda describes as "without respect for due process".[128]
+
+—Jon Lee Anderson, author of Che Guevara: A Revolutionary Life, PBS forum[129]
+
+Although accounts vary, it is estimated that several hundred people were executed nationwide during this time, with Guevara's jurisdictional death total at La Cabaña ranging from 55 to 105.[130] Conflicting views exist of Guevara's attitude towards the executions at La Cabaña. Some exiled opposition biographers report that he relished the rituals of the firing squad, and organized them with gusto, while others relate that Guevara pardoned as many prisoners as he could.[128] All sides acknowledge that Guevara had become a "hardened" man who had no qualms about the death penalty or about summary and collective trials. If the only way to "defend the revolution was to execute its enemies, he would not be swayed by humanitarian or political arguments".[128] In a 5 February 1959, letter to Luis Paredes López in Buenos Aires Guevara states unequivocally: "The executions by firing squads are not only a necessity for the people of Cuba, but also an imposition of the people."[131]
+
+Along with ensuring "revolutionary justice", the other key early platform of Guevara was establishing agrarian land reform. Almost immediately after the success of the revolution, on 27 January 1959, Guevara made one of his most significant speeches where he talked about "the social ideas of the rebel army". During this speech he declared that the main concern of the new Cuban government was "the social justice that land redistribution brings about".[132] A few months later, 17 May 1959, the Agrarian Reform Law, crafted by Guevara, went into effect, limiting the size of all farms to 1,000 acres (400 ha). Any holdings over these limits were expropriated by the government and either redistributed to peasants in 67-acre (270,000 m2) parcels or held as state-run communes.[133] The law also stipulated that foreigners could not own Cuban sugar-plantations.[134]
+
+On 12 June 1959, Castro sent Guevara out on a three-month tour of 14 mostly Bandung Pact countries (Morocco, Sudan, Egypt, Syria, Pakistan, India, Sri Lanka, Burma, Thailand, Indonesia, Japan, Yugoslavia, Greece) and the cities of Singapore and Hong Kong.[135] Sending Guevara away from Havana allowed Castro to appear to distance himself from Guevara and his Marxist sympathies, which troubled both the United States and some of the members of Castro's 26 July Movement.[136] While in Jakarta, Guevara visited Indonesian president Sukarno to discuss the recent revolution of 1945-1949 in Indonesia and to establish trade relations between their two countries. The two men quickly bonded, as Sukarno was attracted to Guevara's energy and his relaxed informal approach; moreover they shared revolutionary leftist aspirations against western imperialism.[137] Guevara next spent 12 days in Japan (15–27 July), participating in negotiations aimed at expanding Cuba's trade relations with that country. During the visit he refused to visit and lay a wreath at Japan's Tomb of the Unknown Soldier commemorating soldiers lost during World War II, remarking that the Japanese "imperialists" had "killed millions of Asians".[138] Instead, Guevara stated that he would visit Hiroshima, where the American military had detonated an atom-bomb 14 years earlier.[138] Despite his denunciation of Imperial Japan, Guevara considered President Truman a "macabre clown" for the bombings,[139] and after visiting Hiroshima and its Peace Memorial Museum he sent back a postcard to Cuba stating, "In order to fight better for peace, one must look at Hiroshima."[140]
+
+Upon Guevara's return to Cuba in September 1959, it became evident that Castro now had more political power. The government had begun land seizures in accordance with the agrarian reform law, but was hedging on compensation offers to landowners, instead offering low-interest "bonds", a step which put the United States on alert. At this point the affected wealthy cattlemen of Camagüey mounted a campaign against the land redistributions and enlisted the newly-disaffected rebel leader Huber Matos, who along with the anti-Communist wing of 26 July Movement, joined them in denouncing "Communist encroachment".[141] During this time Dominican dictator Rafael Trujillo was offering assistance to the "Anti-Communist Legion of the Caribbean" which was training in the Dominican Republic. This multi-national force, composed mostly of Spaniards and Cubans, but also of Croatians, Germans, Greeks, and right-wing mercenaries, was plotting to topple Castro's new regime.[141]
+
+Such threats were heightened when, on 4 March 1960, two massive explosions ripped through the French freighter La Coubre, which was carrying Belgian munitions from the port of Antwerp, and was docked in Havana Harbor. The blasts killed at least 76 people and injured several hundred, with Guevara personally providing first aid to some of the victims. Fidel Castro immediately accused the CIA of "an act of terrorism" and held a state funeral the following day for the victims of the blast.[142] At the memorial service Alberto Korda took the famous photograph of Guevara, now known as Guerrillero Heroico.[143]
+
+Perceived threats prompted Castro to eliminate more "counter-revolutionaries" and to utilize Guevara to drastically increase the speed of land reform. To implement this plan, a new government agency, the National Institute of Agrarian Reform (INRA), was established by the Cuban Government to administer the new Agrarian Reform law. INRA quickly became the most important governing body in the nation, with Guevara serving as its head in his capacity as minister of industries.[134][need quotation to verify] Under Guevara's command, INRA established its own 100,000-person militia, used first to help the government seize control of the expropriated land and supervise its distribution, and later to set up cooperative farms. The land confiscated included 480,000 acres (190,000 ha) owned by United States corporations.[134] Months later, in retaliation, US President Dwight D. Eisenhower sharply reduced United States imports of Cuban sugar (Cuba's main cash crop), which led Guevara on 10 July 1960 to address over 100,000 workers in front of the Presidential Palace at a rally to denounce the "economic aggression" of the United States.[144] Time Magazine reporters who met with Guevara around this time described him as "guid(ing) Cuba with icy calculation, vast competence, high intelligence, and a perceptive sense of humor."[10]
+
+—Urbano (a.k.a. Leonardo Tamayo),fought with Guevara in Cuba and Bolivia[145]
+
+Along with land reform, Guevara stressed the need for national improvement in literacy. Before 1959 the official literacy rate for Cuba was between 60–76%, with educational access in rural areas and a lack of instructors the main determining factors.[146] As a result, the Cuban government at Guevara's behest dubbed 1961 the "year of education" and mobilized over 100,000 volunteers into "literacy brigades", who were then sent out into the countryside to construct schools, train new educators, and teach the predominantly illiterate guajiros (peasants) to read and write.[84][146] Unlike many of Guevara's later economic initiatives, this campaign was "a remarkable success". By the completion of the Cuban Literacy Campaign, 707,212 adults had been taught to read and write, raising the national literacy rate to 96%.[146]
+
+Accompanying literacy, Guevara was also concerned with establishing universal access to higher education. To accomplish this the new regime introduced affirmative action to the universities. While announcing this new commitment, Guevara told the gathered faculty and students at the University of Las Villas that the days when education was "a privilege of the white middle class" had ended. "The University" he said, "must paint itself black, mulatto, worker, and peasant." If it did not, he warned, the people were going to break down its doors "and paint the University the colors they like."[147]
+
+The merit of Marx is that he suddenly produces a qualitative change in the history of social thought. He interprets history, understands its dynamic, predicts the future, but in addition to predicting it (which would satisfy his scientific obligation), he expresses a revolutionary concept: the world must not only be interpreted, it must be transformed. Man ceases to be the slave and tool of his environment and converts himself into the architect of his own destiny.
+
+In September 1960, when Guevara was asked about Cuba's ideology at the First Latin American Congress, he replied, "If I were asked whether our revolution is Communist, I would define it as Marxist. Our revolution has discovered by its methods the paths that Marx pointed out."[149] Consequently, when enacting and advocating Cuban policy, Guevara cited the political philosopher Karl Marx as his ideological inspiration. In defending his political stance, Guevara confidently remarked, "There are truths so evident, so much a part of people's knowledge, that it is now useless to discuss them. One ought to be Marxist with the same naturalness with which one is 'Newtonian' in physics, or 'Pasteurian' in biology."[148] According to Guevara, the "practical revolutionaries" of the Cuban Revolution had the goal of "simply fulfill(ing) laws foreseen by Marx, the scientist."[148] Using Marx's predictions and system of dialectical materialism, Guevara professed that "The laws of Marxism are present in the events of the Cuban Revolution, independently of what its leaders profess or fully know of those laws from a theoretical point of view."[148]
+
+Man truly achieves his full human condition when he produces without being compelled by the physical necessity of selling himself as a commodity.
+
+At this stage, Guevara acquired the additional position of Finance Minister, as well as President of the National Bank.[151] These appointments, combined with his existing position as Minister of Industries, placed Guevara at the zenith of his power, as the "virtual czar" of the Cuban economy.[144] As a consequence of his position at the head of the central bank, it became Guevara's duty to sign the Cuban currency, which per custom bore his signature. Instead of using his full name, he signed the bills solely "Che".[152] It was through this symbolic act, which horrified many in the Cuban financial sector, that Guevara signaled his distaste for money and the class distinctions it brought about.[152] Guevara's long time friend Ricardo Rojo later remarked that "the day he signed Che on the bills, (he) literally knocked the props from under the widespread belief that money was sacred."[153]
+
+In an effort to eliminate social inequalities, Guevara and Cuba's new leadership had moved to swiftly transform the political and economic base of the country through nationalizing factories, banks, and businesses, while attempting to ensure affordable housing, healthcare, and employment for all Cubans.[155] In order for a genuine transformation of consciousness to take root, it was believed that such structural changes had to be accompanied by a conversion in people's social relations and values. Believing that the attitudes in Cuba towards race, women, individualism, and manual labor were the product of the island's outdated past, all individuals were urged to view each other as equals and take on the values of what Guevara termed "el Hombre Nuevo" (the New Man).[155] Guevara hoped his "new man" to be ultimately "selfless and cooperative, obedient and hard working, gender-blind, incorruptible, non-materialistic, and anti-imperialist".[155] To accomplish this, Guevara emphasized the tenets of Marxism–Leninism, and wanted to use the state to emphasize qualities such as egalitarianism and self-sacrifice, at the same time as "unity, equality, and freedom" became the new maxims.[155] Guevara's first desired economic goal of the new man, which coincided with his aversion for wealth condensation and economic inequality, was to see a nationwide elimination of material incentives in favor of moral ones. He negatively viewed capitalism as a "contest among wolves" where "one can only win at the cost of others" and thus desired to see the creation of a "new man and woman".[156] Guevara continually stressed that a socialist economy in itself is not "worth the effort, sacrifice, and risks of war and destruction" if it ends up encouraging "greed and individual ambition at the expense of collective spirit".[157] A primary goal of Guevara's thus became to reform "individual consciousness" and values to produce better workers and citizens.[157] In his view, Cuba's "new man" would be able to overcome the "egotism" and "selfishness" that he loathed and discerned was uniquely characteristic of individuals in capitalist societies.[157] To promote this concept of a "new man", the government also created a series of party-dominated institutions and mechanisms on all levels of society, which included organizations such as labor groups, youth leagues, women's groups, community centers, and houses of culture to promote state-sponsored art, music, and literature. In congruence with this, all educational, mass media, and artistic community based facilities were nationalized and utilized to instill the government's official socialist ideology.[155] In describing this new method of "development", Guevara stated:
+
+There is a great difference between free-enterprise development and revolutionary development. In one of them, wealth is concentrated in the hands of a fortunate few, the friends of the government, the best wheeler-dealers. In the other, wealth is the people's patrimony.[158]
+
+A further integral part of fostering a sense of "unity between the individual and the mass", Guevara believed, was volunteer work and will. To display this, Guevara "led by example", working "endlessly at his ministry job, in construction, and even cutting sugar cane" on his day off.[159] He was known for working 36 hours at a stretch, calling meetings after midnight, and eating on the run.[157] Such behavior was emblematic of Guevara's new program of moral incentives, where each worker was now required to meet a quota and produce a certain quantity of goods. As a replacement for the pay increases abolished by Guevara, workers who exceeded their quota now only received a certificate of commendation, while workers who failed to meet their quotas were given a pay cut.[157] Guevara unapologetically defended his personal philosophy towards motivation and work, stating:
+
+This is not a matter of how many pounds of meat one might be able to eat, or how many times a year someone can go to the beach, or how many ornaments from abroad one might be able to buy with his current salary. What really matters is that the individual feels more complete, with much more internal richness and much more responsibility.[160]
+
+In the face of a loss of commercial connections with Western states, Guevara tried to replace them with closer commercial relationships with Eastern Bloc states, visiting a number of Marxist states and signing trade agreements with them. At the end of 1960 he visited Czechoslovakia, the Soviet Union, North Korea, Hungary and East Germany and signed, for instance, a trade agreement in East Berlin on 17 December 1960.[161] Such agreements helped Cuba's economy to a certain degree but also had the disadvantage of a growing economic dependency on the Eastern Bloc. It was also in East Germany where Guevara met Tamara Bunke (later known as "Tania"), who was assigned as his interpreter, and who joined him years later, and was killed with him in Bolivia.
+
+Whatever the merits or demerits of Guevara's economic principles, his programs were unsuccessful,[162] and accompanied a rapid drop in productivity and a rapid rise in absenteeism.[163] In a meeting with French economist Rene Dumont, Guevara blamed the inadequacy of the Agrarian Reform Law enacted by the Cuban government in 1959, which turned large plantations into farm co-operatives or split up land amongst peasants.[164] In Guevara's opinion, this situation continued to promote a "heightened sense of individual ownership" in which workers could not see the positive social benefits of their labor, leading them to instead seek individual material gain as before.[165] Decades later, Che's former deputy Ernesto Betancourt, the director of Radio Martí, an early ally turned Castro-critic, accused Guevara of being "ignorant of the most elementary economic principles."[166] In reference to the collective failings of Guevara's vision, reporter I. F. Stone who interviewed Guevara twice during this time, remarked that he was "Galahad not Robespierre", while opining that "in a sense he was, like some early saint, taking refuge in the desert. Only there could the purity of the faith be safeguarded from the unregenerate revisionism of human nature".[167]
+
+On 17 April 1961, 1,400 U.S.-trained Cuban exiles invaded Cuba during the Bay of Pigs Invasion. Guevara did not play a key role in the fighting, as one day before the invasion a warship carrying Marines faked an invasion off the West Coast of Pinar del Río and drew forces commanded by Guevara to that region. However, historians give him a share of credit for the victory as he was director of instruction for Cuba's armed forces at the time.[12] Author Tad Szulc in his explanation of the Cuban victory, assigns Guevara partial credit, stating: "The revolutionaries won because Che Guevara, as the head of the Instruction Department of the Revolutionary Armed Forces in charge of the militia training program, had done so well in preparing 200,000 men and women for war."[12] It was also during this deployment that he suffered a bullet grazing to the cheek when his pistol fell out of its holster and accidentally discharged.[168]
+
+In August 1961, during an economic conference of the Organization of American States in Punta del Este, Uruguay, Che Guevara sent a note of "gratitude" to United States President John F. Kennedy through Richard N. Goodwin, Deputy Assistant Secretary of State for Inter-American Affairs. It read "Thanks for Playa Girón (Bay of Pigs). Before the invasion, the revolution was shaky. Now it's stronger than ever."[169] In response to United States Treasury Secretary Douglas Dillon presenting the Alliance for Progress for ratification by the meeting, Guevara antagonistically attacked the United States claim of being a "democracy", stating that such a system was not compatible with "financial oligarchy, discrimination against blacks, and outrages by the Ku Klux Klan".[170] Guevara continued, speaking out against the "persecution" that in his view "drove scientists like Oppenheimer from their posts, deprived the world for years of the marvelous voice of Paul Robeson, and sent the Rosenbergs to their deaths against the protests of a shocked world."[170] Guevara ended his remarks by insinuating that the United States was not interested in real reforms, sardonically quipping that "U.S. experts never talk about agrarian reform; they prefer a safe subject, like a better water supply. In short, they seem to prepare the revolution of the toilets."[171] Nevertheless, Goodwin stated in his memo to President Kennedy following the meeting that Guevara viewed him as someone of the "newer generation"[172] and that Guevara, whom Goodwin alleged sent a message to him the day after the meeting through one of the meeting's Argentine participants whom he described as "Darretta,"[172] also viewed the conversation which the two had as "quite profitable."[172]
+
+Guevara, who was practically the architect of the Soviet–Cuban relationship,[173] then played a key role in bringing to Cuba the Soviet nuclear-armed ballistic missiles that precipitated the Cuban Missile Crisis in October 1962 and brought the world to the brink of nuclear war.[174] A few weeks after the crisis, during an interview with the British communist newspaper the Daily Worker, Guevara was still fuming over the perceived Soviet betrayal and told correspondent Sam Russell that, if the missiles had been under Cuban control, they would have fired them off.[175] While expounding on the incident later, Guevara reiterated that the cause of socialist liberation against global "imperialist aggression" would ultimately have been worth the possibility of "millions of atomic war victims".[176] The missile crisis further convinced Guevara that the world's two superpowers (the United States and the Soviet Union) used Cuba as a pawn in their own global strategies. Afterward, he denounced the Soviets almost as frequently as he denounced the Americans.[177]
+
+In December 1964, Che Guevara had emerged as a "revolutionary statesman of world stature" and thus traveled to New York City as head of the Cuban delegation to speak at the United Nations.[153] On 11 December 1964, during Guevara's hour-long, impassioned address at the UN, he criticized the United Nations' inability to confront the "brutal policy of apartheid" in South Africa, asking "Can the United Nations do nothing to stop this?"[178] Guevara then denounced the United States policy towards their black population, stating:
+
+Those who kill their own children and discriminate daily against them because of the color of their skin; those who let the murderers of blacks remain free, protecting them, and furthermore punishing the black population because they demand their legitimate rights as free men—how can those who do this consider themselves guardians of freedom?[178]
+
+An indignant Guevara ended his speech by reciting the Second Declaration of Havana, decreeing Latin America a "family of 200 million brothers who suffer the same miseries".[178] This "epic", Guevara declared, would be written by the "hungry Indian masses, peasants without land, exploited workers, and progressive masses". To Guevara the conflict was a struggle of masses and ideas, which would be carried forth by those "mistreated and scorned by imperialism" who were previously considered "a weak and submissive flock". With this "flock", Guevara now asserted, "Yankee monopoly capitalism" now terrifyingly saw their "gravediggers".[178] It would be during this "hour of vindication", Guevara pronounced, that the "anonymous mass" would begin to write its own history "with its own blood" and reclaim those "rights that were laughed at by one and all for 500 years". Guevara closed his remarks to the General Assembly by hypothesizing that this "wave of anger" would "sweep the lands of Latin America" and that the labor masses who "turn the wheel of history" were now, for the first time, "awakening from the long, brutalizing sleep to which they had been subjected".[178]
+
+Guevara later learned there had been two failed attempts on his life by Cuban exiles during his stop at the UN complex.[179] The first from Molly Gonzales, who tried to break through barricades upon his arrival with a seven-inch hunting knife, and later during his address by Guillermo Novo, who fired a timer-initiated bazooka from a boat in the East River at the United Nations Headquarters, but missed and was off target. Afterwards Guevara commented on both incidents, stating that "it is better to be killed by a woman with a knife than by a man with a gun", while adding with a languid wave of his cigar that the explosion had "given the whole thing more flavor".[179]
+
+While in New York, Guevara appeared on the CBS Sunday news program Face the Nation,[180] and met with a wide range of people, from United States Senator Eugene McCarthy[181] to associates of Malcolm X. The latter expressed his admiration, declaring Guevara "one of the most revolutionary men in this country right now" while reading a statement from him to a crowd at the Audubon Ballroom.[182]
+
+On 17 December, Guevara left New York for Paris, France, and from there embarked on a three-month world tour that included visits to the People's Republic of China, North Korea, the United Arab Republic, Algeria, Ghana, Guinea, Mali, Dahomey, Congo-Brazzaville and Tanzania, with stops in Ireland and Prague. While in Ireland, Guevara embraced his own Irish heritage, celebrating Saint Patrick's Day in Limerick city.[183] He wrote to his father on this visit, humorously stating "I am in this green Ireland of your ancestors. When they found out, the television [station] came to ask me about the Lynch genealogy, but in case they were horse thieves or something like that, I didn't say much."[184]
+
+During this voyage, he wrote a letter to Carlos Quijano, editor of a Uruguayan weekly, which was later retitled Socialism and Man in Cuba.[156] Outlined in the treatise was Guevara's summons for the creation of a new consciousness, a new status of work, and a new role of the individual. He also laid out the reasoning behind his anti-capitalist sentiments, stating:
+
+The laws of capitalism, blind and invisible to the majority, act upon the individual without his thinking about it. He sees only the vastness of a seemingly infinite horizon before him. That is how it is painted by capitalist propagandists, who purport to draw a lesson from the example of Rockefeller—whether or not it is true—about the possibilities of success. The amount of poverty and suffering required for the emergence of a Rockefeller, and the amount of depravity that the accumulation of a fortune of such magnitude entails, are left out of the picture, and it is not always possible to make the people in general see this.[156]
+
+Guevara ended the essay by declaring that "the true revolutionary is guided by a great feeling of love" and beckoning on all revolutionaries to "strive every day so that this love of living humanity will be transformed into acts that serve as examples", thus becoming "a moving force".[156] The genesis for Guevara's assertions relied on the fact that he believed the example of the Cuban Revolution was "something spiritual that would transcend all borders".[35]
+
+In Algiers, Algeria, on 24 February 1965, Guevara made what turned out to be his last public appearance on the international stage when he delivered a speech at an economic seminar on Afro-Asian solidarity.[185][186] He specified the moral duty of the socialist countries, accusing them of tacit complicity with the exploiting Western countries. He proceeded to outline a number of measures which he said the communist-bloc countries must implement in order to accomplish the defeat of imperialism.[187] Having criticized the Soviet Union (the primary financial backer of Cuba) in such a public manner, he returned to Cuba on 14 March to a solemn reception by Fidel and Raúl Castro, Osvaldo Dorticós and Carlos Rafael Rodríguez at the Havana airport.
+
+As revealed in his last public speech in Algiers, Guevara had come to view the Northern Hemisphere, led by the U.S. in the West and the Soviet Union in the East, as the exploiter of the Southern Hemisphere. He strongly supported Communist North Vietnam in the Vietnam War, and urged the peoples of other developing countries to take up arms and create "many Vietnams".[188] Che's denunciations of the Soviets made him popular among intellectuals and artists of the Western European left who had lost faith in the Soviet Union, while his condemnation of imperialism and call to revolution inspired young radical students in the United States, who were impatient for societal change.[189]
+
+—Helen Yaffe, author of Che Guevara: The Economics of Revolution[190]
+
+In Guevara's private writings from this time (since released), he displays his growing criticism of the Soviet political economy, believing that the Soviets had "forgotten Marx".[190] This led Guevara to denounce a range of Soviet practices including what he saw as their attempt to "air-brush the inherent violence of class struggle integral to the transition from capitalism to socialism", their "dangerous" policy of peaceful co-existence with the United States, their failure to push for a "change in consciousness" towards the idea of work, and their attempt to "liberalize" the socialist economy. Guevara wanted the complete elimination of money, interest, commodity production, the market economy, and "mercantile relationships": all conditions that the Soviets argued would only disappear when world communism was achieved.[190] Disagreeing with this incrementalist approach, Guevara criticized the Soviet Manual of Political Economy, correctly predicting that if USSR did not abolish the law of value (as Guevara desired), it would eventually return to capitalism.[190]
+
+Two weeks after his Algiers speech and his return to Cuba, Guevara dropped out of public life and then vanished altogether.[191] His whereabouts were a great mystery in Cuba, as he was generally regarded as second in power to Castro himself. His disappearance was variously attributed to the failure of the Cuban industrialization scheme he had advocated while minister of industries, to pressure exerted on Castro by Soviet officials who disapproved of Guevara's pro-Chinese Communist stance on the Sino-Soviet split, and to serious differences between Guevara and the pragmatic Castro regarding Cuba's economic development and ideological line.[192] Pressed by international speculation regarding Guevara's fate, Castro stated on 16 June 1965, that the people would be informed when Guevara himself wished to let them know. Still, rumors spread both inside and outside Cuba concerning the missing Guevara's whereabouts.
+
+On 3 October 1965, Castro publicly revealed an undated letter purportedly written to him by Guevara around seven months earlier which was later titled Che Guevara's "farewell letter". In the letter, Guevara reaffirmed his enduring solidarity with the Cuban Revolution but declared his intention to leave Cuba to fight for the revolutionary cause abroad. Additionally, he resigned from all his positions in the Cuban government and communist party, and renounced his honorary Cuban citizenship.[193]
+
+In early 1965, Guevara went to Africa to offer his knowledge and experience as a guerrilla to the ongoing conflict in the Congo. According to Algerian President Ahmed Ben Bella, Guevara thought that Africa was imperialism's weak link and so had enormous revolutionary potential.[194] Egyptian President Gamal Abdel Nasser, who had fraternal relations with Che since his 1959 visit, saw Guevara's plan to fight in Congo as "unwise" and warned that he would become a "Tarzan" figure, doomed to failure.[195] Despite the warning, Guevara traveled to Congo using the alias Ramón Benítez.[196] He led the Cuban operation in support of the Marxist Simba movement, which had emerged from the ongoing Congo crisis. Guevara, his second-in-command Víctor Dreke, and 12 other Cuban expeditionaries arrived in Congo on 24 April 1965, and a contingent of approximately 100 Afro-Cubans joined them soon afterward.[197][198] For a time, they collaborated with guerrilla leader Laurent-Désiré Kabila, who had helped supporters of the overthrown president Patrice Lumumba to lead an unsuccessful revolt months earlier. As an admirer of the late Lumumba, Guevara declared that his "murder should be a lesson for all of us".[199] Guevara, with limited knowledge of Swahili and the local languages, was assigned a teenage interpreter, Freddy Ilanga. Over the course of seven months, Ilanga grew to "admire the hard-working Guevara", who "showed the same respect to black people as he did to whites".[200] Guevara soon became disillusioned with the poor discipline of Kabila's troops and later dismissed him, stating "nothing leads me to believe he is the man of the hour".[201]
+
+As an additional obstacle, white mercenary troops of the Congo National Army, led by Mike Hoare and supported by anti-Castro Cuban pilots and the CIA, thwarted Guevara's movements from his base camp in the mountains near the village of Fizi on Lake Tanganyika in southeast Congo. They were able to monitor his communications and so pre-empted his attacks and interdicted his supply lines. Although Guevara tried to conceal his presence in Congo, the United States government knew his location and activities. The National Security Agency was intercepting all of his incoming and outgoing transmissions via equipment aboard the USNS Private Jose F. Valdez (T-AG-169), a floating listening post that continuously cruised the Indian Ocean off Dar es Salaam for that purpose.[202]
+
+Guevara's aim was to export the revolution by instructing local anti-Mobutu Simba fighters in Marxist ideology and foco theory strategies of guerrilla warfare. In his Congo Diary book, he cites a combination of incompetence, intransigence and infighting among the Congolese rebels as key reasons for the revolt's failure.[203] Later that year, on 20 November 1965, suffering from dysentery and acute asthma, and disheartened after seven months of defeats and inactivity, Guevara left Congo with the six Cuban survivors of his 12-man column. Guevara stated that he had planned to send the wounded back to Cuba and fight in Congo alone until his death, as a revolutionary example. But after being urged by his comrades, and two Cuban emissaries personally sent by Castro, at the last moment he reluctantly agreed to leave Africa. During that day and night, Guevara's forces quietly took down their base camp, burned their huts, and destroyed or threw weapons into Lake Tanganyika that they could not take with them, before crossing the border by boat into Tanzania at night and traveling by land to Dar es Salaam. In speaking about his experience in Congo months later, Guevara concluded that he left rather than fight to the death because: "The human element failed. There is no will to fight. The [rebel] leaders are corrupt. In a word ... there was nothing to do."[204] Guevara also declared that "we can not liberate, all by ourselves, a country that does not want to fight."[205] A few weeks later, he wrote the preface to the diary he kept during the Congo venture, that began: "This is the story of a failure."[206]
+
+Guevara was reluctant to return to Cuba, because Castro had already made public Guevara's "farewell letter"—a letter intended to only be revealed in the case of his death—wherein he severed all ties in order to devote himself to revolution throughout the world.[207] As a result, Guevara spent the next six months living clandestinely at the Cuban embassy in Dar es Salaam and later at a Cuban safehouse in Prague.[208] While in Europe, Guevara made a secret visit to former Argentine president Juan Perón who lived in exile in Francoist Spain where he confided in Perón about his new plan to formulate a communist revolution to bring all of Latin America under socialist control. Perón warned Guevara that his plans for implementing a communist revolution throughout Latin America, starting with Bolivia, would be suicidal and futile, but Guevara's mind was already made up. Later, Perón remarked that Guevara was "an immature utopian... but one of us. I am happy for it to be so because he is giving the Yankees a real headache."[209]
+
+During this time abroad, Guevara compiled his memoirs of the Congo experience and wrote drafts of two more books, one on philosophy and the other on economics. As Guevara prepared for Bolivia, he secretly traveled back to Cuba on 21 July 1966 to visit Castro, as well as to see his wife and to write a last letter to his five children to be read upon his death, which ended with him instructing them:
+
+Above all, always be capable of feeling deeply any injustice committed against anyone, anywhere in the world. This is the most beautiful quality in a revolutionary.[210]
+
+In late 1966, Guevara's location was still not public knowledge, although representatives of Mozambique's independence movement, the FRELIMO, reported that they met with Guevara in late 1966 in Dar es Salaam regarding his offer to aid in their revolutionary project, an offer which they ultimately rejected.[211] In a speech at the 1967 International Workers' Day rally in Havana, the acting minister of the armed forces, Major Juan Almeida, announced that Guevara was "serving the revolution somewhere in Latin America".[citation needed]
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+Before he departed for Bolivia, Guevara altered his appearance by shaving off his beard and much of his hair, also dying it grey so that he was unrecognizable as Che Guevara.[212] On 3 November 1966, Guevara secretly arrived in La Paz on a flight from Montevideo, under the false name Adolfo Mena González, posing as a middle-aged Uruguayan businessman working for the Organization of American States.[213]
+
+Three days after his arrival in Bolivia, Guevara left La Paz for the rural south east region of the country to form his guerrilla army. Guevara's first base camp was located in the montane dry forest in the remote Ñancahuazú region. Training at the camp in the Ñancahuazú valley proved to be hazardous, and little was accomplished in way of building a guerrilla army. The Argentine-born East German operative Haydée Tamara Bunke Bider, better known by her nom de guerre "Tania", had been installed as Che's primary agent in La Paz.[214][215]
+
+Guevara's guerrilla force, numbering about 50 men[216] and operating as the ELN (Ejército de Liberación Nacional de Bolivia, "National Liberation Army of Bolivia"), was well equipped and scored a number of early successes against Bolivian army regulars in the difficult terrain of the mountainous Camiri region during the early months of 1967. As a result of Guevara's units' winning several skirmishes against Bolivian troops in the spring and summer of 1967, the Bolivian government began to overestimate the true size of the guerrilla force.[217] But in August 1967, the Bolivian Army managed to eliminate two guerrilla groups in a violent battle, reportedly killing one of the leaders.[citation needed]
+
+Researchers hypothesize that Guevara's plan for fomenting a revolution in Bolivia failed for an array of reasons:
+
+In addition, Guevara's known preference for confrontation rather than compromise, which had previously surfaced during his guerrilla warfare campaign in Cuba, contributed to his inability to develop successful working relationships with local rebel leaders in Bolivia, just as it had in the Congo.[220] This tendency had existed in Cuba, but had been kept in check by the timely interventions and guidance of Fidel Castro.[221]
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+The end result was that Guevara was unable to attract inhabitants of the local area to join his militia during the eleven months he attempted recruitment. Many of the inhabitants willingly informed the Bolivian authorities and military about the guerrillas and their movements in the area. Near the end of the Bolivian venture, Guevara wrote in his diary that "the peasants do not give us any help, and they are turning into informers."[222]
+
+Félix Rodríguez, a Cuban exile turned CIA Special Activities Division operative, advised Bolivian troops during the hunt for Guevara in Bolivia.[223] In addition, the 2007 documentary My Enemy's Enemy alleges that Nazi war criminal Klaus Barbie advised and possibly helped the CIA orchestrate Guevara's eventual capture.[224]
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+On 7 October 1967, an informant apprised the Bolivian Special Forces of the location of Guevara's guerrilla encampment in the Yuro ravine.[225] On the morning of 8 October, they encircled the area with two battalions numbering 1,800 soldiers and advanced into the ravine triggering a battle where Guevara was wounded and taken prisoner while leading a detachment with Simeón Cuba Sarabia. Che's biographer Jon Lee Anderson reports Bolivian Sergeant Bernardino Huanca's account: that as the Bolivian Rangers approached, a twice-wounded Guevara, his gun rendered useless, threw up his arms in surrender and shouted to the soldiers: "Do not shoot! I am Che Guevara and I am worth more to you alive than dead."[226]
+
+—Philip Agee, CIA agent from 1957–1968, later defected to Cuba [227]
+
+Guevara was tied up and taken to a dilapidated mud schoolhouse in the nearby village of La Higuera on the evening of 8 October. For the next half day, Guevara refused to be interrogated by Bolivian officers and only spoke quietly to Bolivian soldiers. One of those Bolivian soldiers, a helicopter pilot named Jaime Nino de Guzman, describes Che as looking "dreadful". According to Guzman, Guevara was shot through the right calf, his hair was matted with dirt, his clothes were shredded, and his feet were covered in rough leather sheaths. Despite his haggard appearance, he recounts that "Che held his head high, looked everyone straight in the eyes and asked only for something to smoke." De Guzman states that he "took pity" and gave him a small bag of tobacco for his pipe, and that Guevara then smiled and thanked him.[228] Later on the night of 8 October, Guevara—despite having his hands tied—kicked a Bolivian army officer, named Captain Espinosa, against a wall after the officer entered the schoolhouse and tried to snatch Guevara's pipe from his mouth as a souvenir while he was still smoking it.[229] In another instance of defiance, Guevara spat in the face of Bolivian Rear Admiral Ugarteche, who attempted to question Guevara a few hours before his execution.[229]
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+The following morning on 9 October, Guevara asked to see the school teacher of the village, a 22-year-old woman named Julia Cortez. She later stated that she found Guevara to be an "agreeable looking man with a soft and ironic glance" and that during their conversation she found herself "unable to look him in the eye" because his "gaze was unbearable, piercing, and so tranquil".[229] During their short conversation, Guevara pointed out to Cortez the poor condition of the schoolhouse, stating that it was "anti-pedagogical" to expect campesino students to be educated there, while "government officials drive Mercedes cars", and declaring "that's what we are fighting against."[229]
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+Later that morning on 9 October, Bolivian President René Barrientos ordered that Guevara be killed. The order was relayed to the unit holding Guevara by Félix Rodríguez reportedly despite the United States government's desire that Guevara be taken to Panama for further interrogation.[230] The executioner who volunteered to kill Guevara was Mario Terán, a 27-year-old sergeant in the Bolivian army who while half-drunk requested to shoot Guevara because three of his friends from B Company, all with the same first name of "Mario", had been killed in a firefight several days earlier with Guevara's band of guerrillas.[11] To make the bullet wounds appear consistent with the story that the Bolivian government planned to release to the public, Félix Rodríguez ordered Terán not to shoot Guevara in the head, but to aim carefully to make it appear that Guevara had been killed in action during a clash with the Bolivian army.[231] Gary Prado, the Bolivian captain in command of the army company that captured Guevara, said that the reasons Barrientos ordered the immediate execution of Guevara were so there could be no possibility for Guevara to escape from prison, and also so there could be no drama of a public trial where adverse publicity might happen.[232]
+
+About 30 minutes before Guevara was killed, Félix Rodríguez attempted to question him about the whereabouts of other guerrilla fighters who were currently at large, but Guevara continued to remain silent. Rodríguez, assisted by a few Bolivian soldiers, helped Guevara to his feet and took him outside the hut to parade him before other Bolivian soldiers where he posed with Guevara for a photo opportunity where one soldier took a photograph of Rodríguez and other soldiers standing alongside Guevara. Afterwards, Rodríguez told Guevara that he was going to be executed. A little later, Guevara was asked by one of the Bolivian soldiers guarding him if he was thinking about his own immortality. "No," he replied, "I'm thinking about the immortality of the revolution."[233] A few minutes later, Sergeant Terán entered the hut to shoot him, whereupon Guevara reportedly stood up and spoke to Terán what were his last words: "I know you've come to kill me. Shoot, coward! You are only going to kill a man!" Terán hesitated, then pointed his self-loading M2 carbine[234] at Guevara and opened fire, hitting him in the arms and legs.[235] Then, as Guevara writhed on the ground, apparently biting one of his wrists to avoid crying out, Terán fired another burst, fatally wounding him in the chest. Guevara was pronounced dead at 1:10 pm local time according to Rodríguez.[235] In all, Guevara was shot nine times by Terán. This included five times in his legs, once in the right shoulder and arm, and once in the chest and throat.[229]
+
+Months earlier, during his last public declaration to the Tricontinental Conference,[188] Guevara had written his own epitaph, stating: "Wherever death may surprise us, let it be welcome, provided that this our battle cry may have reached some receptive ear and another hand may be extended to wield our weapons."[236]
+
+After his execution, Guevara's body was lashed to the landing skids of a helicopter and flown to nearby Vallegrande, where photographs were taken of him lying on a concrete slab in the laundry room of the Nuestra Señora de Malta.[237] Several witnesses were called to confirm his identity, key amongst them the British journalist Richard Gott, the only witness to have met Guevara when he was alive. Put on display, as hundreds of local residents filed past the body, Guevara's corpse was considered by many to represent a "Christ-like" visage, with some even surreptitiously clipping locks of his hair as divine relics.[238] Such comparisons were further extended when English art critic John Berger, two weeks later upon seeing the post-mortem photographs, observed that they resembled two famous paintings: Rembrandt's The Anatomy Lesson of Dr. Nicolaes Tulp and Andrea Mantegna's Lamentation over the Dead Christ.[239] There were also four correspondents present when Guevara's body arrived in Vallegrande, including Björn Kumm of the Swedish Aftonbladet, who described the scene in a 11 November 1967, exclusive for The New Republic.[240]
+
+A declassified memorandum dated 11 October 1967 to United States President Lyndon B. Johnson from his National Security Advisor Walt Whitman Rostow, called the decision to kill Guevara "stupid" but "understandable from a Bolivian standpoint".[241] After the execution, Rodríguez took several of Guevara's personal items, including a watch which he continued to wear many years later, often showing them to reporters during the ensuing years.[242] Today, some of these belongings, including his flashlight, are on display at the CIA.[243] After a military doctor dismembered his hands, Bolivian army officers transferred Guevara's body to an undisclosed location and refused to reveal whether his remains had been buried or cremated. The hands were sent to Buenos Aires for fingerprint identification. They were later sent to Cuba.[244]
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+On 15 October in Havana, Fidel Castro publicly acknowledged that Guevara was dead and proclaimed three days of public mourning throughout Cuba.[245] On 18 October, Castro addressed a crowd of one million mourners in Havana's Plaza de la Revolución and spoke about Guevara's character as a revolutionary.[246] Fidel Castro closed his impassioned eulogy thus:
+
+If we wish to express what we want the men of future generations to be, we must say: Let them be like Che! If we wish to say how we want our children to be educated, we must say without hesitation: We want them to be educated in Che's spirit! If we want the model of a man, who does not belong to our times but to the future, I say from the depths of my heart that such a model, without a single stain on his conduct, without a single stain on his action, is Che![247]
+
+Also removed when Guevara was captured were his 30,000-word, hand-written diary, a collection of his personal poetry, and a short story he had authored about a young Communist guerrilla who learns to overcome his fears.[248] His diary documented events of the guerrilla campaign in Bolivia,[249] with the first entry on 7 November 1966, shortly after his arrival at the farm in Ñancahuazú, and the last dated 7 October 1967, the day before his capture. The diary tells how the guerrillas were forced to begin operations prematurely because of discovery by the Bolivian Army, explains Guevara's decision to divide the column into two units that were subsequently unable to re-establish contact, and describes their overall unsuccessful venture. It also records the rift between Guevara and the Communist Party of Bolivia that resulted in Guevara having significantly fewer soldiers than originally expected, and shows that Guevara had a great deal of difficulty recruiting from the local populace, partly because the guerrilla group had learned Quechua, unaware that the local language was actually a Tupí–Guaraní language.[250] As the campaign drew to an unexpected close, Guevara became increasingly ill. He suffered from ever-worsening bouts of asthma, and most of his last offensives were carried out in an attempt to obtain medicine.[251] The Bolivian diary was quickly and crudely translated by Ramparts magazine and circulated around the world.[252] There are at least four additional diaries in existence—those of Israel Reyes Zayas (Alias "Braulio"), Harry Villegas Tamayo ("Pombo"), Eliseo Reyes Rodriguez ("Rolando")[214] and Dariel Alarcón Ramírez ("Benigno")[253]—each of which reveals additional aspects of the events.
+
+French intellectual Régis Debray, who was captured in April 1967 while with Guevara in Bolivia, gave an interview from prison in August 1968, in which he enlarged on the circumstances of Guevara's capture. Debray, who had lived with Guevara's band of guerrillas for a short time, said that in his view they were "victims of the forest" and thus "eaten by the jungle".[254] Debray described a destitute situation where Guevara's men suffered malnutrition, lack of water, absence of shoes, and only possessed six blankets for 22 men. Debray recounts that Guevara and the others had been suffering an "illness" which caused their hands and feet to swell into "mounds of flesh" to the point where you could not discern the fingers on their hands. Debray described Guevara as "optimistic about the future of Latin America" despite the futile situation, and remarked that Guevara was "resigned to die in the knowledge that his death would be a sort of renaissance", noting that Guevara perceived death "as a promise of rebirth" and "ritual of renewal".[254]
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+To a certain extent, this belief by Guevara of a metaphorical resurrection came true. While pictures of the dead Guevara were being circulated and the circumstances of his death were being debated, Che's legend began to spread. Demonstrations in protest against his "assassination" occurred throughout the world, and articles, tributes, and poems were written about his life and death.[255] Rallies in support of Guevara were held from "Mexico to Santiago, Algiers to Angola, and Cairo to Calcutta".[256] The population of Budapest and Prague lit candles to honor Guevara's passing; and the picture of a smiling Che appeared in London and Paris.[257] When a few months later riots broke out in Berlin, France, and Chicago, and the unrest spread to the American college campuses, young men and women wore Che Guevara T-shirts and carried his pictures during their protest marches. In the view of military historian Erik Durschmied: "In those heady months of 1968, Che Guevara was not dead. He was very much alive."[258]
+
+In late 1995, the retired Bolivian General Mario Vargas revealed to Jon Lee Anderson, author of Che Guevara: A Revolutionary Life, that Guevara's corpse lay near a Vallegrande airstrip. The result was a multi-national search for the remains, which lasted more than a year. In July 1997 a team of Cuban geologists and Argentine forensic anthropologists discovered the remnants of seven bodies in two mass graves, including one man without hands (as Guevara would have been). Bolivian government officials with the Ministry of Interior later identified the body as Guevara when the excavated teeth "perfectly matched" a plaster mold of Che's teeth made in Cuba prior to his Congolese expedition. The "clincher" then arrived when Argentine forensic anthropologist Alejandro Inchaurregui inspected the inside hidden pocket of a blue jacket dug up next to the handless cadaver and found a small bag of pipe tobacco. Nino de Guzman, the Bolivian helicopter pilot who had given Che a small bag of tobacco, later remarked that he "had serious doubts" at first and "thought the Cubans would just find any old bones and call it Che"; but "after hearing about the tobacco pouch, I have no doubts."[228] On 17 October 1997, Guevara's remains, with those of six of his fellow combatants, were laid to rest with military honors in a specially built mausoleum in the Cuban city of Santa Clara, where he had commanded over the decisive military victory of the Cuban Revolution.[259]
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+In July 2008, the Bolivian government of Evo Morales unveiled Guevara's formerly-sealed diaries composed in two frayed notebooks, along with a logbook and several black-and-white photographs. At this event Bolivia's vice-minister of culture, Pablo Groux, expressed that there were plans to publish photographs of every handwritten page later in the year.[260] Meanwhile, in August 2009 anthropologists working for Bolivia's Justice Ministry discovered and unearthed the bodies of five of Guevara's fellow guerrillas near the Bolivian town of Teoponte.[261]
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+The discovery of Che's remains metonymically activated a series of interlinked associations—rebel, martyr, rogue figure from a picaresque adventure, savior, renegade, extremist—in which there was no fixed divide among them. The current court of opinion places Che on a continuum that teeters between viewing him as a misguided rebel, a coruscatingly brilliant guerrilla philosopher, a poet-warrior jousting at windmills, a brazen warrior who threw down the gauntlet to the bourgeoisie, the object of fervent paeans to his sainthood, or a mass murderer clothed in the guise of an avenging angel whose every action is imbricated in violence—the archetypal Fanatical Terrorist.
+
+Guevara's life and legacy remain contentious. The perceived contradictions of his ethos at various points in his life have created a complex character of duality, one who was "able to wield the pen and submachine gun with equal skill", while prophesying that "the most important revolutionary ambition was to see man liberated from his alienation".[263][264] Guevara's paradoxical standing is further complicated by his array of seemingly diametrically opposed qualities. A secular humanist and sympathetic practitioner of medicine who did not hesitate to shoot his enemies, a celebrated internationalist leader who advocated violence to enforce a utopian philosophy of the collective good, an idealistic intellectual who loved literature but refused to allow dissent, an anti-imperialist Marxist insurgent who was radically willing to forge a poverty-less new world on the apocalyptic ashes of the old one, and finally, an outspoken anti-capitalist whose image has been commoditized. Che's history continues to be rewritten and re-imagined.[265][266] Moreover, sociologist Michael Löwy contends that the many facets of Guevara's life (i.e. doctor and economist, revolutionary and banker, military theoretician and ambassador, deep thinker and political agitator) illuminated the rise of the "Che myth", allowing him to be invariably crystallized in his many metanarrative roles as a "Red Robin Hood, Don Quixote of communism, new Garibaldi, Marxist Saint Just, Cid Campeador of the Wretched of the Earth, Sir Galahad of the beggars ... and Bolshevik devil who haunts the dreams of the rich, kindling braziers of subversion all over the world".[263]
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+As such, various notable individuals have lauded Guevara; for example, Nelson Mandela referred to him as "an inspiration for every human being who loves freedom",[227] while Jean-Paul Sartre described him as "not only an intellectual but also the most complete human being of our age".[267] Others who have expressed their admiration include authors Graham Greene, who remarked that Guevara "represented the idea of gallantry, chivalry, and adventure",[268] and Susan Sontag, who supposed that "[Che's] goal was nothing less than the cause of humanity itself."[269] In the Pan-African community philosopher Frantz Fanon professed Guevara to be "the world symbol of the possibilities of one man",[270] while Black Power leader Stokely Carmichael eulogized that "Che Guevara is not dead, his ideas are with us."[271] Praise has been reflected throughout the political spectrum, with libertarian theorist Murray Rothbard extolling Guevara as a "heroic figure" who "more than any man of our epoch or even of our century, was the living embodiment of the principle of revolution",[272] while journalist Christopher Hitchens reminisced that "[Che's] death meant a lot to me and countless like me at the time, he was a role model, albeit an impossible one for us bourgeois romantics insofar as he went and did what revolutionaries were meant to do—fought and died for his beliefs."[273]
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+Conversely, Jacobo Machover, an exiled opposition author, dismisses all praise of Guevara and portrays him as a callous executioner.[274] Exiled former Cuban prisoners have expressed similar opinions, among them Armando Valladares, who declared Guevara "a man full of hatred" who executed dozens without trial,[275] and Carlos Alberto Montaner, who asserted that Guevara possessed "a Robespierre mentality", wherein cruelty against the revolution's enemies was a virtue.[276] Álvaro Vargas Llosa of The Independent Institute has hypothesized that Guevara's contemporary followers "delude themselves by clinging to a myth", describing Guevara as a "Marxist Puritan" who employed his rigid power to suppress dissent, while also operating as a "cold-blooded killing machine".[166] Llosa also accuses Guevara's "fanatical disposition" as being the linchpin of the "Sovietization" of the Cuban revolution, speculating that he possessed a "total subordination of reality to blind ideological orthodoxy".[166] On a macro-level, Hoover Institution research fellow William Ratliff regards Guevara more as a creation of his historical environment, referring to him as a "fearless" and "head-strong Messiah-like figure", who was the product of a martyr-enamored Latin American culture which "inclined people to seek out and follow paternalistic miracle workers".[277] Ratliff further speculates that the economic conditions in the region suited Guevara's commitment to "bring justice to the downtrodden by crushing centuries-old tyrannies"; describing Latin America as being plagued by what Moisés Naím referred to as the "legendary malignancies" of inequality, poverty, dysfunctional politics and malfunctioning institutions.[277]
+
+In a mixed assessment, British historian Hugh Thomas opined that Guevara was a "brave, sincere and determined man who was also obstinate, narrow, and dogmatic".[279] At the end of his life, according to Thomas, "he seems to have become convinced of the virtues of violence for its own sake", while "his influence over Castro for good or evil" grew after his death, as Fidel took up many of his views.[279] Similarly, the Cuban-American sociologist Samuel Farber lauds Che Guevara as "an honest and committed revolutionary", but also criticizes the fact that "he never embraced socialism in its most democratic essence".[280] Nevertheless, Guevara remains a national hero in Cuba, where his image adorns the 3 peso banknote and school children begin each morning by pledging "We will be like Che."[281][282] In his homeland of Argentina, where high schools bear his name,[283] numerous Che museums dot the country and in 2008 a 12-foot (3.7 m) bronze statue of him was unveiled in the city of his birth, Rosario.[284] Guevara has been sanctified by some Bolivian campesinos[285] as "Saint Ernesto", who pray to him for assistance.[286] In contrast, Guevara remains a hated figure amongst many in the Cuban exile and Cuban-American community of the United States, who view him as "the butcher of La Cabaña".[287] Despite this polarized status, a high-contrast monochrome graphic of Che's face, created in 1968 by Irish artist Jim Fitzpatrick, became a universally merchandized and objectified image,[288][289] found on an endless array of items, including T-shirts, hats, posters, tattoos, and bikinis,[290] contributing to the consumer culture Guevara despised. Yet, he still remains a transcendent figure both in specifically political contexts[291] and as a wide-ranging popular icon of youthful rebellion.[273]
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+Guevara received several honors of state during his life.
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