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--- |
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base_model: stabilityai/stable-diffusion-xl-base-1.0 |
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instance_prompt: pixar-2 |
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tags: |
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- stable-diffusion-xl |
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- stable-diffusion-xl-diffusers |
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- text-to-image |
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- diffusers |
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- lora |
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inference: false |
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datasets: |
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- jbilcke-hf/pixar-2 |
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--- |
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# LoRA DreamBooth - jbilcke-hf/sdxl-pixar-2 |
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## MODEL IS CURRENTLY TRAINING ... |
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Last checkpoint saved: checkpoint-1000 |
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These are LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0 trained on @fffiloni's SD-XL trainer. |
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The weights were trained on the concept prompt: |
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``` |
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pixar-2 |
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``` |
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Use this keyword to trigger your custom model in your prompts. |
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LoRA for the text encoder was enabled: False. |
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Special VAE used for training: madebyollin/sdxl-vae-fp16-fix. |
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## Usage |
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Make sure to upgrade diffusers to >= 0.19.0: |
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``` |
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pip install diffusers --upgrade |
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``` |
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In addition make sure to install transformers, safetensors, accelerate as well as the invisible watermark: |
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``` |
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pip install invisible_watermark transformers accelerate safetensors |
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``` |
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To just use the base model, you can run: |
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```python |
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import torch |
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from diffusers import DiffusionPipeline, AutoencoderKL |
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device = "cuda" if torch.cuda.is_available() else "cpu" |
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vae = AutoencoderKL.from_pretrained('madebyollin/sdxl-vae-fp16-fix', torch_dtype=torch.float16) |
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pipe = DiffusionPipeline.from_pretrained( |
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"stabilityai/stable-diffusion-xl-base-1.0", |
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vae=vae, torch_dtype=torch.float16, variant="fp16", |
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use_safetensors=True |
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) |
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pipe.to(device) |
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# This is where you load your trained weights |
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specific_safetensors = "pytorch_lora_weights.safetensors" |
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lora_scale = 0.9 |
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pipe.load_lora_weights( |
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'jbilcke-hf/sdxl-pixar-2', |
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weight_name = specific_safetensors, |
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# use_auth_token = True |
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) |
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prompt = "A majestic pixar-2 jumping from a big stone at night" |
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image = pipe( |
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prompt=prompt, |
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num_inference_steps=50, |
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cross_attention_kwargs={"scale": lora_scale} |
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).images[0] |
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``` |
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