# Diffusion XL

TL;DR: Enter a prompt or roll the `🎲` and press `Generate`.

## Prompting

Positive and negative prompts are embedded by [Compel](https://github.com/damian0815/compel) for weighting. See [syntax features](https://github.com/damian0815/compel/blob/main/doc/syntax.md) to learn more.

Use `+` or `-` to increase the weight of a token. The weight grows exponentially when chained. For example, `blue+` means 1.1x more attention is given to `blue`, while `blue++` means 1.1^2 more, and so on. The same applies to `-`.

For groups of tokens, wrap them in parentheses and multiply by a float between 0 and 2. For example, `a (birthday cake)1.3 on a table` will increase the weight of both `birthday` and `cake` by 1.3x. This also means the entire scene will be more birthday-like, not just the cake. To counteract this, you can use `-` inside the parentheses on specific tokens, e.g., `a (birthday-- cake)1.3`, to reduce the birthday aspect.

This is the same syntax used in [InvokeAI](https://invoke-ai.github.io/InvokeAI/features/PROMPTS/) and it differs from AUTOMATIC1111:

| Compel      | AUTOMATIC1111 |
| ----------- | ------------- |
| `blue++`    | `((blue))`    |
| `blue--`    | `[[blue]]`    |
| `(blue)1.2` | `(blue:1.2)`  |
| `(blue)0.8` | `(blue:0.8)`  |

### Arrays

Arrays allow you to generate multiple different images from a single prompt. For example, `an adult [[blonde,brunette]] [[man,woman]]` will expand into **4** different prompts. This implementation was inspired by [Fooocus](https://github.com/lllyasviel/Fooocus/pull/1503).

> NB: Make sure to set `Images` to the number of images you want to generate. Otherwise, only the first prompt will be used.

## Models

Each model checkpoint has a different aesthetic:

* [cagliostrolab/animagine-xl-3.1](https://huggingface.co/cagliostrolab/animagine-xl-3.1): anime
* [cyberdelia/CyberRealisticXL](https://huggingface.co/cyberdelia/CyberRealsticXL): photorealistic
* [fluently/Fluently-XL-Final](https://huggingface.co/fluently/Fluently-XL-Final): general purpose
* [segmind/Segmind-Vega](https://huggingface.co/segmind/Segmind-Vega): lightweight general purpose (default)
* [SG161222/RealVisXL_V5.0](https://huggingface.co/SG161222/RealVisXL_V5.0): photorealistic
* [stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0): base

## Styles

[Styles](https://huggingface.co/spaces/adamelliotfields/diffusion/blob/main/data/styles.json) are prompt templates that wrap your positive and negative prompts. They were originally derived from the [twri/sdxl_prompt_styler](https://github.com/twri/sdxl_prompt_styler) Comfy node, but have since been entirely rewritten.

Start by framing a simple subject like `portrait of a young adult woman` or `landscape of a mountain range` and experiment.

## Scale

Rescale up to 4x using [Real-ESRGAN](https://github.com/xinntao/Real-ESRGAN) with weights from [ai-forever](ai-forever/Real-ESRGAN). Necessary for high-resolution images.

## Advanced

### DeepCache

[DeepCache](https://github.com/horseee/DeepCache) caches lower UNet layers and reuses them every `Interval` steps. Trade quality for speed:
* `1`: no caching (default)
* `2`: more quality
* `3`: balanced
* `4`: more speed

### Refiner

Use the [ensemble of expert denoisers](https://research.nvidia.com/labs/dir/eDiff-I/) technique, where the first 80% of timesteps are denoised by the base model and the remaining 80% by the [refiner](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0). Not available with image-to-image pipelines.