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from diffusers import AutoencoderKL, UNet2DConditionModel, StableDiffusionPipeline, StableDiffusionImg2ImgPipeline, DPMSolverMultistepScheduler
import gradio as gr
import torch
from PIL import Image
import os
scheduler = DPMSolverMultistepScheduler(
beta_start=0.00085,
beta_end=0.012,
beta_schedule="scaled_linear",
num_train_timesteps=1000,
trained_betas=None,
predict_epsilon=True,
thresholding=False,
algorithm_type="dpmsolver++",
solver_type="midpoint",
lower_order_final=True,
)
class Model:
def __init__(self, name, path, prefix):
self.name = name
self.path = path
self.prefix = prefix
self.pipe_t2i = None
self.pipe_i2i = None
models = [
Model("Stable-Diffusion-v1.4", "CompVis/stable-diffusion-v1-4", "The 1.4 version of official stable-diffusion"),
Model("Waifu", "hakurei/waifu-diffusion", "anime style"),
]
last_mode = "txt2img"
current_model = models[0]
current_model_path = current_model.path
auth_token = os.getenv("HUGGING_FACE_HUB_TOKEN")
print(f"Is CUDA available: {torch.cuda.is_available()}")
if torch.cuda.is_available():
vae = AutoencoderKL.from_pretrained(current_model.path, subfolder="vae", torch_dtype=torch.float16, use_auth_token=auth_token)
for model in models:
try:
unet = UNet2DConditionModel.from_pretrained(model.path, subfolder="unet", torch_dtype=torch.float16, use_auth_token=auth_token)
model.pipe_t2i = StableDiffusionPipeline.from_pretrained(model.path, unet=unet, vae=vae, torch_dtype=torch.float16, scheduler=scheduler, use_auth_token=auth_token)
model.pipe_i2i = StableDiffusionImg2ImgPipeline.from_pretrained(model.path, unet=unet, vae=vae, torch_dtype=torch.float16, scheduler=scheduler, use_auth_token=auth_token)
except:
models.remove(model)
pipe = models[0].pipe_t2i
pipe = pipe.to("cuda")
else:
vae = AutoencoderKL.from_pretrained(current_model.path, subfolder="vae", use_auth_token=auth_token)
for model in models:
try:
unet = UNet2DConditionModel.from_pretrained(model.path, subfolder="unet", use_auth_token=auth_token)
model.pipe_t2i = StableDiffusionPipeline.from_pretrained(model.path, unet=unet, vae=vae, scheduler=scheduler, use_auth_token=auth_token)
model.pipe_i2i = StableDiffusionImg2ImgPipeline.from_pretrained(model.path, unet=unet, vae=vae, scheduler=scheduler, use_auth_token=auth_token)
except:
models.remove(model)
pipe = models[0].pipe_t2i
pipe = pipe.to("cpu")
device = "GPU 🔥" if torch.cuda.is_available() else "CPU 🥶"
def inference(model_name, prompt, guidance, steps, width=512, height=512, seed=0, img=None, strength=0.5, neg_prompt=""):
global current_model
for model in models:
if model.name == model_name:
current_model = model
model_path = current_model.path
generator = torch.Generator('cuda' if torch.cuda.is_available() else 'cpu').manual_seed(seed) if seed != 0 else None
if img is not None:
return img_to_img(model_path, prompt, neg_prompt, img, strength, guidance, steps, width, height, generator)
else:
return txt_to_img(model_path, prompt, neg_prompt, guidance, steps, width, height, generator)
def txt_to_img(model_path, prompt, neg_prompt, guidance, steps, width, height, generator=None):
global last_mode
global pipe
global current_model_path
if model_path != current_model_path or last_mode != "txt2img":
current_model_path = model_path
pipe.to("cpu")
pipe = current_model.pipe_t2i
if torch.cuda.is_available():
pipe = pipe.to("cuda")
last_mode = "txt2img"
prompt = current_model.prefix + prompt
result = pipe(
prompt,
negative_prompt = neg_prompt,
# num_images_per_prompt=n_images,
num_inference_steps = int(steps),
guidance_scale = guidance,
width = width,
height = height,
generator = generator)
return replace_nsfw_images(result)
def img_to_img(model_path, prompt, neg_prompt, img, strength, guidance, steps, width, height, generator=None):
global last_mode
global pipe
global current_model_path
if model_path != current_model_path or last_mode != "img2img":
current_model_path = model_path
pipe.to("cpu")
pipe = current_model.pipe_i2i
if torch.cuda.is_available():
pipe = pipe.to("cuda")
last_mode = "img2img"
prompt = current_model.prefix + prompt
ratio = min(height / img.height, width / img.width)
img = img.resize((int(img.width * ratio), int(img.height * ratio)), Image.LANCZOS)
result = pipe(
prompt,
negative_prompt = neg_prompt,
# num_images_per_prompt=n_images,
init_image = img,
num_inference_steps = int(steps),
strength = strength,
guidance_scale = guidance,
width = width,
height = height,
generator = generator)
return replace_nsfw_images(result)
def replace_nsfw_images(results):
for i in range(len(results.images)):
if results.nsfw_content_detected[i]:
results.images[i] = Image.open("nsfw.png")
return results.images[0]
css = """
<style>
.finetuned-diffusion-div {
text-align: center;
max-width: 700px;
margin: 0 auto;
font-family: 'IBM Plex Sans', sans-serif;
}
.finetuned-diffusion-div div {
display: inline-flex;
align-items: center;
gap: 0.8rem;
font-size: 1.75rem;
}
.finetuned-diffusion-div div h1 {
font-weight: 900;
margin-top: 15px;
margin-bottom: 15px;
text-align: center;
}
.finetuned-diffusion-div p {
margin-bottom: 10px;
font-size: 94%;
}
.finetuned-diffusion-div p a {
text-decoration: underline;
}
.tabs {
margin-top: 0px;
margin-bottom: 0px;
}
#gallery {
min-height: 20rem;
}
.container {
max-width: 1000px;
margin: auto;
padding-top: 1.5rem;
}
</style>
"""
with gr.Blocks(css=css) as demo:
gr.HTML(
f"""
<div class="finetuned-diffusion-div">
<div>
<h1>Stable-Diffusion with DPM-Solver (fastest sampler for diffusion models) </h1>
</div>
<br>
<p>
❤️ Acknowledgement: Hardware resources of this demo are supported by HuggingFace 🤗 . Many thanks for the help!
</p>
<br>
<p>
This is a demo of sampling by DPM-Solver with two variants of Stable Diffusion models, including <a href="https://huggingface.co/CompVis/stable-diffusion-v1-4">Stable-Diffusion-v1.4</a> and <a href="https://huggingface.co/hakurei/waifu-diffusion">Waifu</a>.
</p>
<br>
<p>
<a href="https://github.com/LuChengTHU/dpm-solver">DPM-Solver</a> (Neurips 2022 Oral) is a fast high-order solver customized for diffusion ODEs, which can generate high-quality samples by diffusion models within only 10-25 steps. DPM-Solver has an analytical formulation and is very easy to use for all types of Gaussian diffusion models, and includes <a href="https://arxiv.org/abs/2010.02502">DDIM</a> as a first-order special case.
</p>
<p>
We use <a href="https://github.com/huggingface/diffusers">Diffusers</a> 🧨 to implement this demo, which currently supports the multistep DPM-Solver scheduler. For more details of DPM-Solver with Diffusers, check <a href="https://github.com/huggingface/diffusers/pull/1132">this pull request</a>.
</p>
<br>
<p>
Currently, the default sampler of stable-diffusion is <a href="https://arxiv.org/abs/2202.09778">PNDM</a>, which needs 50 steps to generate high-quality samples. However, DPM-Solver can generate high-quality samples within only <span style="font-weight: bold;">20-25</span> steps, and for some samples even within <span style="font-weight: bold;">10-15</span> steps.
</p>
<br>
<p>
Running on <b>{device}</b>
</p>
</div>
"""
)
with gr.Row():
with gr.Column(scale=55):
with gr.Group():
model_name = gr.Dropdown(label="Model", choices=[m.name for m in models], value=current_model.name)
with gr.Row():
prompt = gr.Textbox(label="Prompt", show_label=False, max_lines=2,placeholder="Enter prompt. Style applied automatically").style(container=False)
generate = gr.Button(value="Generate").style(rounded=(False, True, True, False))
image_out = gr.Image(height=512)
# gallery = gr.Gallery(
# label="Generated images", show_label=False, elem_id="gallery"
# ).style(grid=[1], height="auto")
with gr.Column(scale=45):
with gr.Tab("Options"):
with gr.Group():
neg_prompt = gr.Textbox(label="Negative prompt", placeholder="What to exclude from the image")
# n_images = gr.Slider(label="Images", value=1, minimum=1, maximum=4, step=1)
with gr.Row():
guidance = gr.Slider(label="Guidance scale", value=7.5, maximum=15)
steps = gr.Slider(label="Steps", value=25, minimum=2, maximum=100, step=1)
with gr.Row():
width = gr.Slider(label="Width", value=512, minimum=64, maximum=1024, step=8)
height = gr.Slider(label="Height", value=512, minimum=64, maximum=1024, step=8)
seed = gr.Slider(0, 2147483647, label='Seed (0 = random)', value=0, step=1)
with gr.Tab("Image to image"):
with gr.Group():
image = gr.Image(label="Image", height=256, tool="editor", type="pil")
strength = gr.Slider(label="Transformation strength", minimum=0, maximum=1, step=0.01, value=0.5)
# model_name.change(lambda x: gr.update(visible = x == models[0].name), inputs=model_name, outputs=custom_model_group)
inputs = [model_name, prompt, guidance, steps, width, height, seed, image, strength, neg_prompt]
prompt.submit(inference, inputs=inputs, outputs=image_out)
generate.click(inference, inputs=inputs, outputs=image_out)
gr.Markdown('''
Stable-diffusion Models by [CompVis](https://huggingface.co/CompVis) and [stabilityai](https://huggingface.co/stabilityai), Waifu-diffusion models by [@hakurei](https://huggingface.co/hakurei). Most of the code of this demo are copied from [@anzorq's fintuned-diffusion](https://huggingface.co/spaces/anzorq/finetuned_diffusion/tree/main) ❤️<br>
Space by [Cheng Lu](https://github.com/LuChengTHU). [![Twitter Follow](https://img.shields.io/twitter/follow/ChengLu05671218?label=%40ChengLu&style=social)](https://twitter.com/ChengLu05671218)
![visitors](https://visitor-badge.glitch.me/badge?page_id=LuChengTHU.dpmsolver_sdm)
''')
demo.queue(concurrency_count=1)
demo.launch(debug=False, share=False) |