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# Copyright 2023 The HuggingFace Team. All rights reserved. | |
# | |
# Licensed under the Apache License, Version 2.0 (the "License"); | |
# you may not use this file except in compliance with the License. | |
# You may obtain a copy of the License at | |
# | |
# http://www.apache.org/licenses/LICENSE-2.0 | |
# | |
# Unless required by applicable law or agreed to in writing, software | |
# distributed under the License is distributed on an "AS IS" BASIS, | |
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. | |
# See the License for the specific language governing permissions and | |
# limitations under the License. | |
# This model implementation is heavily inspired by https://github.com/haofanwang/ControlNet-for-Diffusers/ | |
import inspect | |
import warnings | |
from typing import Any, Callable, Dict, List, Optional, Tuple, Union | |
import numpy as np | |
import PIL.Image | |
import torch | |
import torch.nn.functional as F | |
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer | |
from diffusers.image_processor import VaeImageProcessor | |
from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin | |
from diffusers.models import AutoencoderKL, ControlNetModel, UNet2DConditionModel | |
from diffusers.schedulers import KarrasDiffusionSchedulers | |
from diffusers.utils import ( | |
is_accelerate_available, | |
is_accelerate_version, | |
logging, | |
replace_example_docstring, | |
) | |
from diffusers.utils.torch_utils import randn_tensor,is_compiled_module | |
from diffusers.pipelines.pipeline_utils import DiffusionPipeline | |
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput | |
from diffusers.pipelines.stable_diffusion.safety_checker import StableDiffusionSafetyChecker | |
from diffusers.pipelines.controlnet import MultiControlNetModel | |
logger = logging.get_logger(__name__) # pylint: disable=invalid-name | |
EXAMPLE_DOC_STRING = """ | |
Examples: | |
```py | |
>>> # !pip install transformers accelerate | |
>>> from diffusers import StableDiffusionControlNetInpaintPipeline, ControlNetModel, DDIMScheduler | |
>>> from diffusers.utils import load_image | |
>>> import numpy as np | |
>>> import torch | |
>>> init_image = load_image( | |
... "https://huggingface.co/datasets/diffusers/test-arrays/resolve/main/stable_diffusion_inpaint/boy.png" | |
... ) | |
>>> init_image = init_image.resize((512, 512)) | |
>>> generator = torch.Generator(device="cpu").manual_seed(1) | |
>>> mask_image = load_image( | |
... "https://huggingface.co/datasets/diffusers/test-arrays/resolve/main/stable_diffusion_inpaint/boy_mask.png" | |
... ) | |
>>> mask_image = mask_image.resize((512, 512)) | |
>>> def make_inpaint_condition(image, image_mask): | |
... image = np.array(image.convert("RGB")).astype(np.float32) / 255.0 | |
... image_mask = np.array(image_mask.convert("L")).astype(np.float32) / 255.0 | |
... assert image.shape[0:1] == image_mask.shape[0:1], "image and image_mask must have the same image size" | |
... image[image_mask > 0.5] = -1.0 # set as masked pixel | |
... image = np.expand_dims(image, 0).transpose(0, 3, 1, 2) | |
... image = torch.from_numpy(image) | |
... return image | |
>>> control_image = make_inpaint_condition(init_image, mask_image) | |
>>> controlnet = ControlNetModel.from_pretrained( | |
... "lllyasviel/control_v11p_sd15_inpaint", torch_dtype=torch.float16 | |
... ) | |
>>> pipe = StableDiffusionControlNetInpaintPipeline.from_pretrained( | |
... "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16 | |
... ) | |
>>> pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config) | |
>>> pipe.enable_model_cpu_offload() | |
>>> # generate image | |
>>> image = pipe( | |
... "a handsome man with ray-ban sunglasses", | |
... num_inference_steps=20, | |
... generator=generator, | |
... eta=1.0, | |
... image=init_image, | |
... mask_image=mask_image, | |
... control_image=control_image, | |
... ).images[0] | |
``` | |
""" | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint.prepare_mask_and_masked_image | |
def prepare_mask_and_masked_image(image, mask, height, width, return_image=False): | |
""" | |
Prepares a pair (image, mask) to be consumed by the Stable Diffusion pipeline. This means that those inputs will be | |
converted to ``torch.Tensor`` with shapes ``batch x channels x height x width`` where ``channels`` is ``3`` for the | |
``image`` and ``1`` for the ``mask``. | |
The ``image`` will be converted to ``torch.float32`` and normalized to be in ``[-1, 1]``. The ``mask`` will be | |
binarized (``mask > 0.5``) and cast to ``torch.float32`` too. | |
Args: | |
image (Union[np.array, PIL.Image, torch.Tensor]): The image to inpaint. | |
It can be a ``PIL.Image``, or a ``height x width x 3`` ``np.array`` or a ``channels x height x width`` | |
``torch.Tensor`` or a ``batch x channels x height x width`` ``torch.Tensor``. | |
mask (_type_): The mask to apply to the image, i.e. regions to inpaint. | |
It can be a ``PIL.Image``, or a ``height x width`` ``np.array`` or a ``1 x height x width`` | |
``torch.Tensor`` or a ``batch x 1 x height x width`` ``torch.Tensor``. | |
Raises: | |
ValueError: ``torch.Tensor`` images should be in the ``[-1, 1]`` range. ValueError: ``torch.Tensor`` mask | |
should be in the ``[0, 1]`` range. ValueError: ``mask`` and ``image`` should have the same spatial dimensions. | |
TypeError: ``mask`` is a ``torch.Tensor`` but ``image`` is not | |
(ot the other way around). | |
Returns: | |
tuple[torch.Tensor]: The pair (mask, masked_image) as ``torch.Tensor`` with 4 | |
dimensions: ``batch x channels x height x width``. | |
""" | |
if image is None: | |
raise ValueError("`image` input cannot be undefined.") | |
if mask is None: | |
raise ValueError("`mask_image` input cannot be undefined.") | |
if isinstance(image, torch.Tensor): | |
if not isinstance(mask, torch.Tensor): | |
raise TypeError(f"`image` is a torch.Tensor but `mask` (type: {type(mask)} is not") | |
# Batch single image | |
if image.ndim == 3: | |
assert image.shape[0] == 3, "Image outside a batch should be of shape (3, H, W)" | |
image = image.unsqueeze(0) | |
# Batch and add channel dim for single mask | |
if mask.ndim == 2: | |
mask = mask.unsqueeze(0).unsqueeze(0) | |
# Batch single mask or add channel dim | |
if mask.ndim == 3: | |
# Single batched mask, no channel dim or single mask not batched but channel dim | |
if mask.shape[0] == 1: | |
mask = mask.unsqueeze(0) | |
# Batched masks no channel dim | |
else: | |
mask = mask.unsqueeze(1) | |
assert image.ndim == 4 and mask.ndim == 4, "Image and Mask must have 4 dimensions" | |
assert image.shape[-2:] == mask.shape[-2:], "Image and Mask must have the same spatial dimensions" | |
assert image.shape[0] == mask.shape[0], "Image and Mask must have the same batch size" | |
# Check image is in [-1, 1] | |
if image.min() < -1 or image.max() > 1: | |
raise ValueError("Image should be in [-1, 1] range") | |
# Check mask is in [0, 1] | |
if mask.min() < 0 or mask.max() > 1: | |
raise ValueError("Mask should be in [0, 1] range") | |
# Binarize mask | |
mask[mask < 0.5] = 0 | |
mask[mask >= 0.5] = 1 | |
# Image as float32 | |
image = image.to(dtype=torch.float32) | |
elif isinstance(mask, torch.Tensor): | |
raise TypeError(f"`mask` is a torch.Tensor but `image` (type: {type(image)} is not") | |
else: | |
# preprocess image | |
if isinstance(image, (PIL.Image.Image, np.ndarray)): | |
image = [image] | |
if isinstance(image, list) and isinstance(image[0], PIL.Image.Image): | |
# resize all images w.r.t passed height an width | |
image = [i.resize((width, height), resample=PIL.Image.LANCZOS) for i in image] | |
image = [np.array(i.convert("RGB"))[None, :] for i in image] | |
image = np.concatenate(image, axis=0) | |
elif isinstance(image, list) and isinstance(image[0], np.ndarray): | |
image = np.concatenate([i[None, :] for i in image], axis=0) | |
image = image.transpose(0, 3, 1, 2) | |
image = torch.from_numpy(image).to(dtype=torch.float32) / 127.5 - 1.0 | |
# preprocess mask | |
if isinstance(mask, (PIL.Image.Image, np.ndarray)): | |
mask = [mask] | |
if isinstance(mask, list) and isinstance(mask[0], PIL.Image.Image): | |
mask = [i.resize((width, height), resample=PIL.Image.LANCZOS) for i in mask] | |
mask = np.concatenate([np.array(m.convert("L"))[None, None, :] for m in mask], axis=0) | |
mask = mask.astype(np.float32) / 255.0 | |
elif isinstance(mask, list) and isinstance(mask[0], np.ndarray): | |
mask = np.concatenate([m[None, None, :] for m in mask], axis=0) | |
mask[mask < 0.5] = 0 | |
mask[mask >= 0.5] = 1 | |
mask = torch.from_numpy(mask) | |
masked_image = image * (mask < 0.5) | |
# n.b. ensure backwards compatibility as old function does not return image | |
if return_image: | |
return mask, masked_image, image | |
return mask, masked_image | |
class StableDiffusionControlNetInpaintPipeline( | |
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin | |
): | |
r""" | |
Pipeline for text-to-image generation using Stable Diffusion with ControlNet guidance. | |
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the | |
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.) | |
In addition the pipeline inherits the following loading methods: | |
- *Textual-Inversion*: [`loaders.TextualInversionLoaderMixin.load_textual_inversion`] | |
<Tip> | |
This pipeline can be used both with checkpoints that have been specifically fine-tuned for inpainting, such as | |
[runwayml/stable-diffusion-inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting) | |
as well as default text-to-image stable diffusion checkpoints, such as | |
[runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5). | |
Default text-to-image stable diffusion checkpoints might be preferable for controlnets that have been fine-tuned on | |
those, such as [lllyasviel/control_v11p_sd15_inpaint](https://huggingface.co/lllyasviel/control_v11p_sd15_inpaint). | |
</Tip> | |
Args: | |
vae ([`AutoencoderKL`]): | |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations. | |
text_encoder ([`CLIPTextModel`]): | |
Frozen text-encoder. Stable Diffusion uses the text portion of | |
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically | |
the [clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14) variant. | |
tokenizer (`CLIPTokenizer`): | |
Tokenizer of class | |
[CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer). | |
unet ([`UNet2DConditionModel`]): Conditional U-Net architecture to denoise the encoded image latents. | |
controlnet ([`ControlNetModel`] or `List[ControlNetModel]`): | |
Provides additional conditioning to the unet during the denoising process. If you set multiple ControlNets | |
as a list, the outputs from each ControlNet are added together to create one combined additional | |
conditioning. | |
scheduler ([`SchedulerMixin`]): | |
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of | |
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`]. | |
safety_checker ([`StableDiffusionSafetyChecker`]): | |
Classification module that estimates whether generated images could be considered offensive or harmful. | |
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details. | |
feature_extractor ([`CLIPImageProcessor`]): | |
Model that extracts features from generated images to be used as inputs for the `safety_checker`. | |
""" | |
_optional_components = ["safety_checker", "feature_extractor"] | |
def __init__( | |
self, | |
vae: AutoencoderKL, | |
text_encoder: CLIPTextModel, | |
tokenizer: CLIPTokenizer, | |
unet: UNet2DConditionModel, | |
controlnet: Union[ControlNetModel, List[ControlNetModel], Tuple[ControlNetModel], MultiControlNetModel], | |
scheduler: KarrasDiffusionSchedulers, | |
safety_checker: StableDiffusionSafetyChecker, | |
feature_extractor: CLIPImageProcessor, | |
requires_safety_checker: bool = True, | |
): | |
super().__init__() | |
if safety_checker is None and requires_safety_checker: | |
logger.warning( | |
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure" | |
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered" | |
" results in services or applications open to the public. Both the diffusers team and Hugging Face" | |
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling" | |
" it only for use-cases that involve analyzing network behavior or auditing its results. For more" | |
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ." | |
) | |
if safety_checker is not None and feature_extractor is None: | |
raise ValueError( | |
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety" | |
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead." | |
) | |
if isinstance(controlnet, (list, tuple)): | |
controlnet = MultiControlNetModel(controlnet) | |
self.register_modules( | |
vae=vae, | |
text_encoder=text_encoder, | |
tokenizer=tokenizer, | |
unet=unet, | |
controlnet=controlnet, | |
scheduler=scheduler, | |
safety_checker=safety_checker, | |
feature_extractor=feature_extractor, | |
) | |
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1) | |
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor) | |
self.control_image_processor = VaeImageProcessor( | |
vae_scale_factor=self.vae_scale_factor, do_convert_rgb=True, do_normalize=False | |
) | |
self.register_to_config(requires_safety_checker=requires_safety_checker) | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_slicing | |
def enable_vae_slicing(self): | |
r""" | |
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to | |
compute decoding in several steps. This is useful to save some memory and allow larger batch sizes. | |
""" | |
self.vae.enable_slicing() | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_slicing | |
def disable_vae_slicing(self): | |
r""" | |
Disable sliced VAE decoding. If `enable_vae_slicing` was previously enabled, this method will go back to | |
computing decoding in one step. | |
""" | |
self.vae.disable_slicing() | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_tiling | |
def enable_vae_tiling(self): | |
r""" | |
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to | |
compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow | |
processing larger images. | |
""" | |
self.vae.enable_tiling() | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_tiling | |
def disable_vae_tiling(self): | |
r""" | |
Disable tiled VAE decoding. If `enable_vae_tiling` was previously enabled, this method will go back to | |
computing decoding in one step. | |
""" | |
self.vae.disable_tiling() | |
def enable_model_cpu_offload(self, gpu_id=0): | |
r""" | |
Offloads all models to CPU using accelerate, reducing memory usage with a low impact on performance. Compared | |
to `enable_sequential_cpu_offload`, this method moves one whole model at a time to the GPU when its `forward` | |
method is called, and the model remains in GPU until the next model runs. Memory savings are lower than with | |
`enable_sequential_cpu_offload`, but performance is much better due to the iterative execution of the `unet`. | |
""" | |
if is_accelerate_available() and is_accelerate_version(">=", "0.17.0.dev0"): | |
from accelerate import cpu_offload_with_hook | |
else: | |
raise ImportError("`enable_model_cpu_offload` requires `accelerate v0.17.0` or higher.") | |
device = torch.device(f"cuda:{gpu_id}") | |
hook = None | |
for cpu_offloaded_model in [self.text_encoder, self.unet, self.vae]: | |
_, hook = cpu_offload_with_hook(cpu_offloaded_model, device, prev_module_hook=hook) | |
if self.safety_checker is not None: | |
# the safety checker can offload the vae again | |
_, hook = cpu_offload_with_hook(self.safety_checker, device, prev_module_hook=hook) | |
# control net hook has be manually offloaded as it alternates with unet | |
cpu_offload_with_hook(self.controlnet, device) | |
# We'll offload the last model manually. | |
self.final_offload_hook = hook | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._encode_prompt | |
def _encode_prompt( | |
self, | |
promptA, | |
promptB, | |
t, | |
device, | |
num_images_per_prompt, | |
do_classifier_free_guidance, | |
negative_promptA=None, | |
negative_promptB=None, | |
t_nag = None, | |
prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
lora_scale: Optional[float] = None, | |
): | |
r""" | |
Encodes the prompt into text encoder hidden states. | |
Args: | |
prompt (`str` or `List[str]`, *optional*): | |
prompt to be encoded | |
device: (`torch.device`): | |
torch device | |
num_images_per_prompt (`int`): | |
number of images that should be generated per prompt | |
do_classifier_free_guidance (`bool`): | |
whether to use classifier free guidance or not | |
negative_prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation. If not defined, one has to pass | |
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is | |
less than `1`). | |
prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not | |
provided, text embeddings will be generated from `prompt` input argument. | |
negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input | |
argument. | |
lora_scale (`float`, *optional*): | |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded. | |
""" | |
# set lora scale so that monkey patched LoRA | |
# function of text encoder can correctly access it | |
if lora_scale is not None and isinstance(self, LoraLoaderMixin): | |
self._lora_scale = lora_scale | |
prompt = promptA | |
negative_prompt = negative_promptA | |
if promptA is not None and isinstance(promptA, str): | |
batch_size = 1 | |
elif promptA is not None and isinstance(promptA, list): | |
batch_size = len(promptA) | |
else: | |
batch_size = prompt_embeds.shape[0] | |
if prompt_embeds is None: | |
# textual inversion: procecss multi-vector tokens if necessary | |
if isinstance(self, TextualInversionLoaderMixin): | |
promptA = self.maybe_convert_prompt(promptA, self.tokenizer) | |
text_inputsA = self.tokenizer( | |
promptA, | |
padding="max_length", | |
max_length=self.tokenizer.model_max_length, | |
truncation=True, | |
return_tensors="pt", | |
) | |
text_inputsB = self.tokenizer( | |
promptB, | |
padding="max_length", | |
max_length=self.tokenizer.model_max_length, | |
truncation=True, | |
return_tensors="pt", | |
) | |
text_input_idsA = text_inputsA.input_ids | |
text_input_idsB = text_inputsB.input_ids | |
untruncated_ids = self.tokenizer(promptA, padding="longest", return_tensors="pt").input_ids | |
if untruncated_ids.shape[-1] >= text_input_idsA.shape[-1] and not torch.equal( | |
text_input_idsA, untruncated_ids | |
): | |
removed_text = self.tokenizer.batch_decode( | |
untruncated_ids[:, self.tokenizer.model_max_length - 1 : -1] | |
) | |
logger.warning( | |
"The following part of your input was truncated because CLIP can only handle sequences up to" | |
f" {self.tokenizer.model_max_length} tokens: {removed_text}" | |
) | |
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask: | |
attention_mask = text_inputsA.attention_mask.to(device) | |
else: | |
attention_mask = None | |
# print("text_input_idsA: ",text_input_idsA) | |
# print("text_input_idsB: ",text_input_idsB) | |
# print('t: ',t) | |
prompt_embedsA = self.text_encoder( | |
text_input_idsA.to(device), | |
attention_mask=attention_mask, | |
) | |
prompt_embedsA = prompt_embedsA[0] | |
prompt_embedsB = self.text_encoder( | |
text_input_idsB.to(device), | |
attention_mask=attention_mask, | |
) | |
prompt_embedsB = prompt_embedsB[0] | |
prompt_embeds = prompt_embedsA*(t)+(1-t)*prompt_embedsB | |
# print("prompt_embeds: ",prompt_embeds) | |
if self.text_encoder is not None: | |
prompt_embeds_dtype = self.text_encoder.dtype | |
elif self.unet is not None: | |
prompt_embeds_dtype = self.unet.dtype | |
else: | |
prompt_embeds_dtype = prompt_embeds.dtype | |
prompt_embeds = prompt_embeds.to(dtype=prompt_embeds_dtype, device=device) | |
bs_embed, seq_len, _ = prompt_embeds.shape | |
# duplicate text embeddings for each generation per prompt, using mps friendly method | |
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1) | |
prompt_embeds = prompt_embeds.view(bs_embed * num_images_per_prompt, seq_len, -1) | |
# get unconditional embeddings for classifier free guidance | |
if do_classifier_free_guidance and negative_prompt_embeds is None: | |
uncond_tokensA: List[str] | |
uncond_tokensB: List[str] | |
if negative_prompt is None: | |
uncond_tokensA = [""] * batch_size | |
uncond_tokensB = [""] * batch_size | |
elif prompt is not None and type(prompt) is not type(negative_prompt): | |
raise TypeError( | |
f"`negative_prompt` should be the same type to `prompt`, but got {type(negative_prompt)} !=" | |
f" {type(prompt)}." | |
) | |
elif isinstance(negative_prompt, str): | |
uncond_tokensA = [negative_promptA] | |
uncond_tokensB = [negative_promptB] | |
elif batch_size != len(negative_prompt): | |
raise ValueError( | |
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:" | |
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches" | |
" the batch size of `prompt`." | |
) | |
else: | |
uncond_tokensA = negative_promptA | |
uncond_tokensB = negative_promptB | |
# textual inversion: procecss multi-vector tokens if necessary | |
if isinstance(self, TextualInversionLoaderMixin): | |
uncond_tokensA = self.maybe_convert_prompt(uncond_tokensA, self.tokenizer) | |
uncond_tokensB = self.maybe_convert_prompt(uncond_tokensB, self.tokenizer) | |
max_length = prompt_embeds.shape[1] | |
uncond_inputA = self.tokenizer( | |
uncond_tokensA, | |
padding="max_length", | |
max_length=max_length, | |
truncation=True, | |
return_tensors="pt", | |
) | |
uncond_inputB = self.tokenizer( | |
uncond_tokensB, | |
padding="max_length", | |
max_length=max_length, | |
truncation=True, | |
return_tensors="pt", | |
) | |
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask: | |
attention_mask = uncond_inputA.attention_mask.to(device) | |
else: | |
attention_mask = None | |
negative_prompt_embedsA = self.text_encoder( | |
uncond_inputA.input_ids.to(device), | |
attention_mask=attention_mask, | |
) | |
negative_prompt_embedsB = self.text_encoder( | |
uncond_inputB.input_ids.to(device), | |
attention_mask=attention_mask, | |
) | |
negative_prompt_embeds = negative_prompt_embedsA[0]*(t_nag)+(1-t_nag)*negative_prompt_embedsB[0] | |
# negative_prompt_embeds = negative_prompt_embeds[0] | |
if do_classifier_free_guidance: | |
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method | |
seq_len = negative_prompt_embeds.shape[1] | |
negative_prompt_embeds = negative_prompt_embeds.to(dtype=prompt_embeds_dtype, device=device) | |
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1) | |
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1) | |
# For classifier free guidance, we need to do two forward passes. | |
# Here we concatenate the unconditional and text embeddings into a single batch | |
# to avoid doing two forward passes | |
# print("prompt_embeds: ",prompt_embeds) | |
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds]) | |
return prompt_embeds | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.run_safety_checker | |
def run_safety_checker(self, image, device, dtype): | |
if self.safety_checker is None: | |
has_nsfw_concept = None | |
else: | |
if torch.is_tensor(image): | |
feature_extractor_input = self.image_processor.postprocess(image, output_type="pil") | |
else: | |
feature_extractor_input = self.image_processor.numpy_to_pil(image) | |
safety_checker_input = self.feature_extractor(feature_extractor_input, return_tensors="pt").to(device) | |
image, has_nsfw_concept = self.safety_checker( | |
images=image, clip_input=safety_checker_input.pixel_values.to(dtype) | |
) | |
return image, has_nsfw_concept | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.decode_latents | |
def decode_latents(self, latents): | |
warnings.warn( | |
"The decode_latents method is deprecated and will be removed in a future version. Please" | |
" use VaeImageProcessor instead", | |
FutureWarning, | |
) | |
latents = 1 / self.vae.config.scaling_factor * latents | |
image = self.vae.decode(latents, return_dict=False)[0] | |
image = (image / 2 + 0.5).clamp(0, 1) | |
# we always cast to float32 as this does not cause significant overhead and is compatible with bfloat16 | |
image = image.cpu().permute(0, 2, 3, 1).float().numpy() | |
return image | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs | |
def prepare_extra_step_kwargs(self, generator, eta): | |
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature | |
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers. | |
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502 | |
# and should be between [0, 1] | |
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys()) | |
extra_step_kwargs = {} | |
if accepts_eta: | |
extra_step_kwargs["eta"] = eta | |
# check if the scheduler accepts generator | |
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys()) | |
if accepts_generator: | |
extra_step_kwargs["generator"] = generator | |
return extra_step_kwargs | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_img2img.StableDiffusionImg2ImgPipeline.get_timesteps | |
def get_timesteps(self, num_inference_steps, strength, device): | |
# get the original timestep using init_timestep | |
init_timestep = min(int(num_inference_steps * strength), num_inference_steps) | |
t_start = max(num_inference_steps - init_timestep, 0) | |
timesteps = self.scheduler.timesteps[t_start * self.scheduler.order :] | |
return timesteps, num_inference_steps - t_start | |
def check_inputs( | |
self, | |
prompt, | |
image, | |
height, | |
width, | |
callback_steps, | |
negative_prompt=None, | |
prompt_embeds=None, | |
negative_prompt_embeds=None, | |
controlnet_conditioning_scale=1.0, | |
control_guidance_start=0.0, | |
control_guidance_end=1.0, | |
): | |
if height % 8 != 0 or width % 8 != 0: | |
raise ValueError(f"`height` and `width` have to be divisible by 8 but are {height} and {width}.") | |
if (callback_steps is None) or ( | |
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0) | |
): | |
raise ValueError( | |
f"`callback_steps` has to be a positive integer but is {callback_steps} of type" | |
f" {type(callback_steps)}." | |
) | |
if prompt is not None and prompt_embeds is not None: | |
raise ValueError( | |
f"Cannot forward both `prompt`: {prompt} and `prompt_embeds`: {prompt_embeds}. Please make sure to" | |
" only forward one of the two." | |
) | |
elif prompt is None and prompt_embeds is None: | |
raise ValueError( | |
"Provide either `prompt` or `prompt_embeds`. Cannot leave both `prompt` and `prompt_embeds` undefined." | |
) | |
elif prompt is not None and (not isinstance(prompt, str) and not isinstance(prompt, list)): | |
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}") | |
if negative_prompt is not None and negative_prompt_embeds is not None: | |
raise ValueError( | |
f"Cannot forward both `negative_prompt`: {negative_prompt} and `negative_prompt_embeds`:" | |
f" {negative_prompt_embeds}. Please make sure to only forward one of the two." | |
) | |
if prompt_embeds is not None and negative_prompt_embeds is not None: | |
if prompt_embeds.shape != negative_prompt_embeds.shape: | |
raise ValueError( | |
"`prompt_embeds` and `negative_prompt_embeds` must have the same shape when passed directly, but" | |
f" got: `prompt_embeds` {prompt_embeds.shape} != `negative_prompt_embeds`" | |
f" {negative_prompt_embeds.shape}." | |
) | |
# `prompt` needs more sophisticated handling when there are multiple | |
# conditionings. | |
if isinstance(self.controlnet, MultiControlNetModel): | |
if isinstance(prompt, list): | |
logger.warning( | |
f"You have {len(self.controlnet.nets)} ControlNets and you have passed {len(prompt)}" | |
" prompts. The conditionings will be fixed across the prompts." | |
) | |
# Check `image` | |
is_compiled = hasattr(F, "scaled_dot_product_attention") and isinstance( | |
self.controlnet, torch._dynamo.eval_frame.OptimizedModule | |
) | |
if ( | |
isinstance(self.controlnet, ControlNetModel) | |
or is_compiled | |
and isinstance(self.controlnet._orig_mod, ControlNetModel) | |
): | |
self.check_image(image, prompt, prompt_embeds) | |
elif ( | |
isinstance(self.controlnet, MultiControlNetModel) | |
or is_compiled | |
and isinstance(self.controlnet._orig_mod, MultiControlNetModel) | |
): | |
if not isinstance(image, list): | |
raise TypeError("For multiple controlnets: `image` must be type `list`") | |
# When `image` is a nested list: | |
# (e.g. [[canny_image_1, pose_image_1], [canny_image_2, pose_image_2]]) | |
elif any(isinstance(i, list) for i in image): | |
raise ValueError("A single batch of multiple conditionings are supported at the moment.") | |
elif len(image) != len(self.controlnet.nets): | |
raise ValueError( | |
f"For multiple controlnets: `image` must have the same length as the number of controlnets, but got {len(image)} images and {len(self.controlnet.nets)} ControlNets." | |
) | |
for image_ in image: | |
self.check_image(image_, prompt, prompt_embeds) | |
else: | |
assert False | |
# Check `controlnet_conditioning_scale` | |
if ( | |
isinstance(self.controlnet, ControlNetModel) | |
or is_compiled | |
and isinstance(self.controlnet._orig_mod, ControlNetModel) | |
): | |
if not isinstance(controlnet_conditioning_scale, float): | |
raise TypeError("For single controlnet: `controlnet_conditioning_scale` must be type `float`.") | |
elif ( | |
isinstance(self.controlnet, MultiControlNetModel) | |
or is_compiled | |
and isinstance(self.controlnet._orig_mod, MultiControlNetModel) | |
): | |
if isinstance(controlnet_conditioning_scale, list): | |
if any(isinstance(i, list) for i in controlnet_conditioning_scale): | |
raise ValueError("A single batch of multiple conditionings are supported at the moment.") | |
elif isinstance(controlnet_conditioning_scale, list) and len(controlnet_conditioning_scale) != len( | |
self.controlnet.nets | |
): | |
raise ValueError( | |
"For multiple controlnets: When `controlnet_conditioning_scale` is specified as `list`, it must have" | |
" the same length as the number of controlnets" | |
) | |
else: | |
assert False | |
if len(control_guidance_start) != len(control_guidance_end): | |
raise ValueError( | |
f"`control_guidance_start` has {len(control_guidance_start)} elements, but `control_guidance_end` has {len(control_guidance_end)} elements. Make sure to provide the same number of elements to each list." | |
) | |
if isinstance(self.controlnet, MultiControlNetModel): | |
if len(control_guidance_start) != len(self.controlnet.nets): | |
raise ValueError( | |
f"`control_guidance_start`: {control_guidance_start} has {len(control_guidance_start)} elements but there are {len(self.controlnet.nets)} controlnets available. Make sure to provide {len(self.controlnet.nets)}." | |
) | |
for start, end in zip(control_guidance_start, control_guidance_end): | |
if start >= end: | |
raise ValueError( | |
f"control guidance start: {start} cannot be larger or equal to control guidance end: {end}." | |
) | |
if start < 0.0: | |
raise ValueError(f"control guidance start: {start} can't be smaller than 0.") | |
if end > 1.0: | |
raise ValueError(f"control guidance end: {end} can't be larger than 1.0.") | |
# Copied from diffusers.pipelines.controlnet.pipeline_controlnet.StableDiffusionControlNetPipeline.check_image | |
def check_image(self, image, prompt, prompt_embeds): | |
image_is_pil = isinstance(image, PIL.Image.Image) | |
image_is_tensor = isinstance(image, torch.Tensor) | |
image_is_np = isinstance(image, np.ndarray) | |
image_is_pil_list = isinstance(image, list) and isinstance(image[0], PIL.Image.Image) | |
image_is_tensor_list = isinstance(image, list) and isinstance(image[0], torch.Tensor) | |
image_is_np_list = isinstance(image, list) and isinstance(image[0], np.ndarray) | |
if ( | |
not image_is_pil | |
and not image_is_tensor | |
and not image_is_np | |
and not image_is_pil_list | |
and not image_is_tensor_list | |
and not image_is_np_list | |
): | |
raise TypeError( | |
f"image must be passed and be one of PIL image, numpy array, torch tensor, list of PIL images, list of numpy arrays or list of torch tensors, but is {type(image)}" | |
) | |
if image_is_pil: | |
image_batch_size = 1 | |
else: | |
image_batch_size = len(image) | |
if prompt is not None and isinstance(prompt, str): | |
prompt_batch_size = 1 | |
elif prompt is not None and isinstance(prompt, list): | |
prompt_batch_size = len(prompt) | |
elif prompt_embeds is not None: | |
prompt_batch_size = prompt_embeds.shape[0] | |
if image_batch_size != 1 and image_batch_size != prompt_batch_size: | |
raise ValueError( | |
f"If image batch size is not 1, image batch size must be same as prompt batch size. image batch size: {image_batch_size}, prompt batch size: {prompt_batch_size}" | |
) | |
# Copied from diffusers.pipelines.controlnet.pipeline_controlnet.StableDiffusionControlNetPipeline.prepare_image | |
def prepare_control_image( | |
self, | |
image, | |
width, | |
height, | |
batch_size, | |
num_images_per_prompt, | |
device, | |
dtype, | |
do_classifier_free_guidance=False, | |
guess_mode=False, | |
): | |
image = self.control_image_processor.preprocess(image, height=height, width=width).to(dtype=torch.float32) | |
image_batch_size = image.shape[0] | |
if image_batch_size == 1: | |
repeat_by = batch_size | |
else: | |
# image batch size is the same as prompt batch size | |
repeat_by = num_images_per_prompt | |
image = image.repeat_interleave(repeat_by, dim=0) | |
image = image.to(device=device, dtype=dtype) | |
if do_classifier_free_guidance and not guess_mode: | |
image = torch.cat([image] * 2) | |
return image | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint.StableDiffusionInpaintPipeline.prepare_latents | |
def prepare_latents( | |
self, | |
batch_size, | |
num_channels_latents, | |
height, | |
width, | |
dtype, | |
device, | |
generator, | |
latents=None, | |
image=None, | |
timestep=None, | |
is_strength_max=True, | |
return_noise=False, | |
return_image_latents=False, | |
): | |
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor) | |
if isinstance(generator, list) and len(generator) != batch_size: | |
raise ValueError( | |
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch" | |
f" size of {batch_size}. Make sure the batch size matches the length of the generators." | |
) | |
if (image is None or timestep is None) and not is_strength_max: | |
raise ValueError( | |
"Since strength < 1. initial latents are to be initialised as a combination of Image + Noise." | |
"However, either the image or the noise timestep has not been provided." | |
) | |
if return_image_latents or (latents is None and not is_strength_max): | |
image = image.to(device=device, dtype=dtype) | |
image_latents = self._encode_vae_image(image=image, generator=generator) | |
if latents is None: | |
noise = randn_tensor(shape, generator=generator, device=device, dtype=dtype) | |
# if strength is 1. then initialise the latents to noise, else initial to image + noise | |
latents = noise if is_strength_max else self.scheduler.add_noise(image_latents, noise, timestep) | |
# if pure noise then scale the initial latents by the Scheduler's init sigma | |
latents = latents * self.scheduler.init_noise_sigma if is_strength_max else latents | |
else: | |
noise = latents.to(device) | |
latents = noise * self.scheduler.init_noise_sigma | |
outputs = (latents,) | |
if return_noise: | |
outputs += (noise,) | |
if return_image_latents: | |
outputs += (image_latents,) | |
return outputs | |
def _default_height_width(self, height, width, image): | |
# NOTE: It is possible that a list of images have different | |
# dimensions for each image, so just checking the first image | |
# is not _exactly_ correct, but it is simple. | |
while isinstance(image, list): | |
image = image[0] | |
if height is None: | |
if isinstance(image, PIL.Image.Image): | |
height = image.height | |
elif isinstance(image, torch.Tensor): | |
height = image.shape[2] | |
height = (height // 8) * 8 # round down to nearest multiple of 8 | |
if width is None: | |
if isinstance(image, PIL.Image.Image): | |
width = image.width | |
elif isinstance(image, torch.Tensor): | |
width = image.shape[3] | |
width = (width // 8) * 8 # round down to nearest multiple of 8 | |
return height, width | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint.StableDiffusionInpaintPipeline.prepare_mask_latents | |
def prepare_mask_latents( | |
self, mask, masked_image, batch_size, height, width, dtype, device, generator, do_classifier_free_guidance | |
): | |
# resize the mask to latents shape as we concatenate the mask to the latents | |
# we do that before converting to dtype to avoid breaking in case we're using cpu_offload | |
# and half precision | |
mask = torch.nn.functional.interpolate( | |
mask, size=(height // self.vae_scale_factor, width // self.vae_scale_factor) | |
) | |
mask = mask.to(device=device, dtype=dtype) | |
masked_image = masked_image.to(device=device, dtype=dtype) | |
masked_image_latents = self._encode_vae_image(masked_image, generator=generator) | |
# duplicate mask and masked_image_latents for each generation per prompt, using mps friendly method | |
if mask.shape[0] < batch_size: | |
if not batch_size % mask.shape[0] == 0: | |
raise ValueError( | |
"The passed mask and the required batch size don't match. Masks are supposed to be duplicated to" | |
f" a total batch size of {batch_size}, but {mask.shape[0]} masks were passed. Make sure the number" | |
" of masks that you pass is divisible by the total requested batch size." | |
) | |
mask = mask.repeat(batch_size // mask.shape[0], 1, 1, 1) | |
if masked_image_latents.shape[0] < batch_size: | |
if not batch_size % masked_image_latents.shape[0] == 0: | |
raise ValueError( | |
"The passed images and the required batch size don't match. Images are supposed to be duplicated" | |
f" to a total batch size of {batch_size}, but {masked_image_latents.shape[0]} images were passed." | |
" Make sure the number of images that you pass is divisible by the total requested batch size." | |
) | |
masked_image_latents = masked_image_latents.repeat(batch_size // masked_image_latents.shape[0], 1, 1, 1) | |
mask = torch.cat([mask] * 2) if do_classifier_free_guidance else mask | |
masked_image_latents = ( | |
torch.cat([masked_image_latents] * 2) if do_classifier_free_guidance else masked_image_latents | |
) | |
# aligning device to prevent device errors when concating it with the latent model input | |
masked_image_latents = masked_image_latents.to(device=device, dtype=dtype) | |
return mask, masked_image_latents | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint.StableDiffusionInpaintPipeline._encode_vae_image | |
def _encode_vae_image(self, image: torch.Tensor, generator: torch.Generator): | |
if isinstance(generator, list): | |
image_latents = [ | |
self.vae.encode(image[i : i + 1]).latent_dist.sample(generator=generator[i]) | |
for i in range(image.shape[0]) | |
] | |
image_latents = torch.cat(image_latents, dim=0) | |
else: | |
image_latents = self.vae.encode(image).latent_dist.sample(generator=generator) | |
image_latents = self.vae.config.scaling_factor * image_latents | |
return image_latents | |
def predict_woControl( | |
self, | |
promptA: Union[str, List[str]] = None, | |
promptB: Union[str, List[str]] = None, | |
image: Union[torch.FloatTensor, PIL.Image.Image] = None, | |
mask_image: Union[torch.FloatTensor, PIL.Image.Image] = None, | |
height: Optional[int] = None, | |
width: Optional[int] = None, | |
strength: float = 1.0, | |
tradoff: float = 1.0, | |
tradoff_nag: float = 1.0, | |
num_inference_steps: int = 50, | |
guidance_scale: float = 7.5, | |
negative_promptA: Optional[Union[str, List[str]]] = None, | |
negative_promptB: Optional[Union[str, List[str]]] = None, | |
num_images_per_prompt: Optional[int] = 1, | |
eta: float = 0.0, | |
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None, | |
latents: Optional[torch.FloatTensor] = None, | |
prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
output_type: Optional[str] = "pil", | |
return_dict: bool = True, | |
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None, | |
callback_steps: int = 1, | |
cross_attention_kwargs: Optional[Dict[str, Any]] = None, | |
task_class: Union[torch.Tensor, float, int] = None, | |
): | |
r""" | |
The call function to the pipeline for generation. | |
Args: | |
prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts to guide image generation. If not defined, you need to pass `prompt_embeds`. | |
image (`PIL.Image.Image`): | |
`Image` or tensor representing an image batch to be inpainted (which parts of the image to be masked | |
out with `mask_image` and repainted according to `prompt`). | |
mask_image (`PIL.Image.Image`): | |
`Image` or tensor representing an image batch to mask `image`. White pixels in the mask are repainted | |
while black pixels are preserved. If `mask_image` is a PIL image, it is converted to a single channel | |
(luminance) before use. If it's a tensor, it should contain one color channel (L) instead of 3, so the | |
expected shape would be `(B, H, W, 1)`. | |
height (`int`, *optional*, defaults to `self.unet.config.sample_size * self.vae_scale_factor`): | |
The height in pixels of the generated image. | |
width (`int`, *optional*, defaults to `self.unet.config.sample_size * self.vae_scale_factor`): | |
The width in pixels of the generated image. | |
strength (`float`, *optional*, defaults to 1.0): | |
Indicates extent to transform the reference `image`. Must be between 0 and 1. `image` is used as a | |
starting point and more noise is added the higher the `strength`. The number of denoising steps depends | |
on the amount of noise initially added. When `strength` is 1, added noise is maximum and the denoising | |
process runs for the full number of iterations specified in `num_inference_steps`. A value of 1 | |
essentially ignores `image`. | |
num_inference_steps (`int`, *optional*, defaults to 50): | |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the | |
expense of slower inference. This parameter is modulated by `strength`. | |
guidance_scale (`float`, *optional*, defaults to 7.5): | |
A higher guidance scale value encourages the model to generate images closely linked to the text | |
`prompt` at the expense of lower image quality. Guidance scale is enabled when `guidance_scale > 1`. | |
negative_prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to | |
pass `negative_prompt_embeds` instead. Ignored when not using guidance (`guidance_scale < 1`). | |
num_images_per_prompt (`int`, *optional*, defaults to 1): | |
The number of images to generate per prompt. | |
eta (`float`, *optional*, defaults to 0.0): | |
Corresponds to parameter eta (η) from the [DDIM](https://arxiv.org/abs/2010.02502) paper. Only applies | |
to the [`~schedulers.DDIMScheduler`], and is ignored in other schedulers. | |
generator (`torch.Generator` or `List[torch.Generator]`, *optional*): | |
A [`torch.Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make | |
generation deterministic. | |
latents (`torch.FloatTensor`, *optional*): | |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image | |
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents | |
tensor is generated by sampling using the supplied random `generator`. | |
prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not | |
provided, text embeddings are generated from the `prompt` input argument. | |
negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If | |
not provided, `negative_prompt_embeds` are generated from the `negative_prompt` input argument. | |
output_type (`str`, *optional*, defaults to `"pil"`): | |
The output format of the generated image. Choose between `PIL.Image` or `np.array`. | |
return_dict (`bool`, *optional*, defaults to `True`): | |
Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] instead of a | |
plain tuple. | |
callback (`Callable`, *optional*): | |
A function that calls every `callback_steps` steps during inference. The function is called with the | |
following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`. | |
callback_steps (`int`, *optional*, defaults to 1): | |
The frequency at which the `callback` function is called. If not specified, the callback is called at | |
every step. | |
cross_attention_kwargs (`dict`, *optional*): | |
A kwargs dictionary that if specified is passed along to the [`AttentionProcessor`] as defined in | |
[`self.processor`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py). | |
Examples: | |
```py | |
>>> import PIL | |
>>> import requests | |
>>> import torch | |
>>> from io import BytesIO | |
>>> from diffusers import StableDiffusionInpaintPipeline | |
>>> def download_image(url): | |
... response = requests.get(url) | |
... return PIL.Image.open(BytesIO(response.content)).convert("RGB") | |
>>> img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" | |
>>> mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" | |
>>> init_image = download_image(img_url).resize((512, 512)) | |
>>> mask_image = download_image(mask_url).resize((512, 512)) | |
>>> pipe = StableDiffusionInpaintPipeline.from_pretrained( | |
... "runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16 | |
... ) | |
>>> pipe = pipe.to("cuda") | |
>>> prompt = "Face of a yellow cat, high resolution, sitting on a park bench" | |
>>> image = pipe(prompt=prompt, image=init_image, mask_image=mask_image).images[0] | |
``` | |
Returns: | |
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] or `tuple`: | |
If `return_dict` is `True`, [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] is returned, | |
otherwise a `tuple` is returned where the first element is a list with the generated images and the | |
second element is a list of `bool`s indicating whether the corresponding generated image contains | |
"not-safe-for-work" (nsfw) content. | |
""" | |
# 0. Default height and width to unet | |
height = height or self.unet.config.sample_size * self.vae_scale_factor | |
width = width or self.unet.config.sample_size * self.vae_scale_factor | |
prompt = promptA | |
negative_prompt = negative_promptA | |
# 1. Check inputs | |
self.check_inputs( | |
prompt, | |
height, | |
width, | |
strength, | |
callback_steps, | |
negative_prompt, | |
prompt_embeds, | |
negative_prompt_embeds, | |
) | |
# 2. Define call parameters | |
if prompt is not None and isinstance(prompt, str): | |
batch_size = 1 | |
elif prompt is not None and isinstance(prompt, list): | |
batch_size = len(prompt) | |
else: | |
batch_size = prompt_embeds.shape[0] | |
device = self._execution_device | |
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2) | |
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1` | |
# corresponds to doing no classifier free guidance. | |
do_classifier_free_guidance = guidance_scale > 1.0 | |
# 3. Encode input prompt | |
text_encoder_lora_scale = ( | |
cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None | |
) | |
prompt_embeds = self._encode_prompt( | |
promptA, | |
promptB, | |
tradoff, | |
device, | |
num_images_per_prompt, | |
do_classifier_free_guidance, | |
negative_promptA, | |
negative_promptB, | |
tradoff_nag, | |
prompt_embeds=prompt_embeds, | |
negative_prompt_embeds=negative_prompt_embeds, | |
lora_scale=text_encoder_lora_scale, | |
) | |
# 4. set timesteps | |
self.scheduler.set_timesteps(num_inference_steps, device=device) | |
timesteps, num_inference_steps = self.get_timesteps( | |
num_inference_steps=num_inference_steps, strength=strength, device=device | |
) | |
# check that number of inference steps is not < 1 - as this doesn't make sense | |
if num_inference_steps < 1: | |
raise ValueError( | |
f"After adjusting the num_inference_steps by strength parameter: {strength}, the number of pipeline" | |
f"steps is {num_inference_steps} which is < 1 and not appropriate for this pipeline." | |
) | |
# at which timestep to set the initial noise (n.b. 50% if strength is 0.5) | |
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt) | |
# create a boolean to check if the strength is set to 1. if so then initialise the latents with pure noise | |
is_strength_max = strength == 1.0 | |
# 5. Preprocess mask and image | |
mask, masked_image, init_image = prepare_mask_and_masked_image( | |
image, mask_image, height, width, return_image=True | |
) | |
mask_condition = mask.clone() | |
# 6. Prepare latent variables | |
num_channels_latents = self.vae.config.latent_channels | |
num_channels_unet = self.unet.config.in_channels | |
return_image_latents = num_channels_unet == 4 | |
latents_outputs = self.prepare_latents( | |
batch_size * num_images_per_prompt, | |
num_channels_latents, | |
height, | |
width, | |
prompt_embeds.dtype, | |
device, | |
generator, | |
latents, | |
image=init_image, | |
timestep=latent_timestep, | |
is_strength_max=is_strength_max, | |
return_noise=True, | |
return_image_latents=return_image_latents, | |
) | |
if return_image_latents: | |
latents, noise, image_latents = latents_outputs | |
else: | |
latents, noise = latents_outputs | |
# 7. Prepare mask latent variables | |
mask, masked_image_latents = self.prepare_mask_latents( | |
mask, | |
masked_image, | |
batch_size * num_images_per_prompt, | |
height, | |
width, | |
prompt_embeds.dtype, | |
device, | |
generator, | |
do_classifier_free_guidance, | |
) | |
# 8. Check that sizes of mask, masked image and latents match | |
if num_channels_unet == 9: | |
# default case for runwayml/stable-diffusion-inpainting | |
num_channels_mask = mask.shape[1] | |
num_channels_masked_image = masked_image_latents.shape[1] | |
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels: | |
raise ValueError( | |
f"Incorrect configuration settings! The config of `pipeline.unet`: {self.unet.config} expects" | |
f" {self.unet.config.in_channels} but received `num_channels_latents`: {num_channels_latents} +" | |
f" `num_channels_mask`: {num_channels_mask} + `num_channels_masked_image`: {num_channels_masked_image}" | |
f" = {num_channels_latents+num_channels_masked_image+num_channels_mask}. Please verify the config of" | |
" `pipeline.unet` or your `mask_image` or `image` input." | |
) | |
elif num_channels_unet != 4: | |
raise ValueError( | |
f"The unet {self.unet.__class__} should have either 4 or 9 input channels, not {self.unet.config.in_channels}." | |
) | |
# 9. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline | |
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta) | |
# 10. Denoising loop | |
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order | |
with self.progress_bar(total=num_inference_steps) as progress_bar: | |
for i, t in enumerate(timesteps): | |
# expand the latents if we are doing classifier free guidance | |
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents | |
# concat latents, mask, masked_image_latents in the channel dimension | |
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t) | |
if num_channels_unet == 9: | |
latent_model_input = torch.cat([latent_model_input, mask, masked_image_latents], dim=1) | |
# predict the noise residual | |
if task_class is not None: | |
noise_pred = self.unet( | |
sample = latent_model_input, | |
timestep = t, | |
encoder_hidden_states=prompt_embeds, | |
cross_attention_kwargs=cross_attention_kwargs, | |
return_dict=False, | |
task_class = task_class, | |
)[0] | |
else: | |
noise_pred = self.unet( | |
latent_model_input, | |
t, | |
encoder_hidden_states=prompt_embeds, | |
cross_attention_kwargs=cross_attention_kwargs, | |
return_dict=False, | |
)[0] | |
# perform guidance | |
if do_classifier_free_guidance: | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
# compute the previous noisy sample x_t -> x_t-1 | |
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0] | |
if num_channels_unet == 4: | |
init_latents_proper = image_latents[:1] | |
init_mask = mask[:1] | |
if i < len(timesteps) - 1: | |
noise_timestep = timesteps[i + 1] | |
init_latents_proper = self.scheduler.add_noise( | |
init_latents_proper, noise, torch.tensor([noise_timestep]) | |
) | |
latents = (1 - init_mask) * init_latents_proper + init_mask * latents | |
# call the callback, if provided | |
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0): | |
progress_bar.update() | |
if callback is not None and i % callback_steps == 0: | |
callback(i, t, latents) | |
if not output_type == "latent": | |
condition_kwargs = {} | |
if isinstance(self.vae, AsymmetricAutoencoderKL): | |
init_image = init_image.to(device=device, dtype=masked_image_latents.dtype) | |
init_image_condition = init_image.clone() | |
init_image = self._encode_vae_image(init_image, generator=generator) | |
mask_condition = mask_condition.to(device=device, dtype=masked_image_latents.dtype) | |
condition_kwargs = {"image": init_image_condition, "mask": mask_condition} | |
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, **condition_kwargs)[0] | |
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype) | |
else: | |
image = latents | |
has_nsfw_concept = None | |
if has_nsfw_concept is None: | |
do_denormalize = [True] * image.shape[0] | |
else: | |
do_denormalize = [not has_nsfw for has_nsfw in has_nsfw_concept] | |
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize) | |
# Offload last model to CPU | |
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None: | |
self.final_offload_hook.offload() | |
if not return_dict: | |
return (image, has_nsfw_concept) | |
return StableDiffusionPipelineOutput(images=image, nsfw_content_detected=has_nsfw_concept) | |
def __call__( | |
self, | |
promptA: Union[str, List[str]] = None, | |
promptB: Union[str, List[str]] = None, | |
image: Union[torch.Tensor, PIL.Image.Image] = None, | |
mask_image: Union[torch.Tensor, PIL.Image.Image] = None, | |
control_image: Union[ | |
torch.FloatTensor, | |
PIL.Image.Image, | |
np.ndarray, | |
List[torch.FloatTensor], | |
List[PIL.Image.Image], | |
List[np.ndarray], | |
] = None, | |
height: Optional[int] = None, | |
width: Optional[int] = None, | |
strength: float = 1.0, | |
tradoff: float = 1.0, | |
tradoff_nag: float = 1.0, | |
num_inference_steps: int = 50, | |
guidance_scale: float = 7.5, | |
negative_promptA: Optional[Union[str, List[str]]] = None, | |
negative_promptB: Optional[Union[str, List[str]]] = None, | |
num_images_per_prompt: Optional[int] = 1, | |
eta: float = 0.0, | |
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None, | |
latents: Optional[torch.FloatTensor] = None, | |
prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
output_type: Optional[str] = "pil", | |
return_dict: bool = True, | |
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None, | |
callback_steps: int = 1, | |
cross_attention_kwargs: Optional[Dict[str, Any]] = None, | |
controlnet_conditioning_scale: Union[float, List[float]] = 0.5, | |
guess_mode: bool = False, | |
control_guidance_start: Union[float, List[float]] = 0.0, | |
control_guidance_end: Union[float, List[float]] = 1.0, | |
): | |
r""" | |
Function invoked when calling the pipeline for generation. | |
Args: | |
prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`. | |
instead. | |
image (`torch.FloatTensor`, `PIL.Image.Image`, `List[torch.FloatTensor]`, `List[PIL.Image.Image]`, | |
`List[List[torch.FloatTensor]]`, or `List[List[PIL.Image.Image]]`): | |
The ControlNet input condition. ControlNet uses this input condition to generate guidance to Unet. If | |
the type is specified as `Torch.FloatTensor`, it is passed to ControlNet as is. `PIL.Image.Image` can | |
also be accepted as an image. The dimensions of the output image defaults to `image`'s dimensions. If | |
height and/or width are passed, `image` is resized according to them. If multiple ControlNets are | |
specified in init, images must be passed as a list such that each element of the list can be correctly | |
batched for input to a single controlnet. | |
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor): | |
The height in pixels of the generated image. | |
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor): | |
The width in pixels of the generated image. | |
strength (`float`, *optional*, defaults to 1.): | |
Conceptually, indicates how much to transform the masked portion of the reference `image`. Must be | |
between 0 and 1. `image` will be used as a starting point, adding more noise to it the larger the | |
`strength`. The number of denoising steps depends on the amount of noise initially added. When | |
`strength` is 1, added noise will be maximum and the denoising process will run for the full number of | |
iterations specified in `num_inference_steps`. A value of 1, therefore, essentially ignores the masked | |
portion of the reference `image`. | |
num_inference_steps (`int`, *optional*, defaults to 50): | |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the | |
expense of slower inference. | |
guidance_scale (`float`, *optional*, defaults to 7.5): | |
Guidance scale as defined in [Classifier-Free Diffusion Guidance](https://arxiv.org/abs/2207.12598). | |
`guidance_scale` is defined as `w` of equation 2. of [Imagen | |
Paper](https://arxiv.org/pdf/2205.11487.pdf). Guidance scale is enabled by setting `guidance_scale > | |
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`, | |
usually at the expense of lower image quality. | |
negative_prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation. If not defined, one has to pass | |
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is | |
less than `1`). | |
num_images_per_prompt (`int`, *optional*, defaults to 1): | |
The number of images to generate per prompt. | |
eta (`float`, *optional*, defaults to 0.0): | |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to | |
[`schedulers.DDIMScheduler`], will be ignored for others. | |
generator (`torch.Generator` or `List[torch.Generator]`, *optional*): | |
One or a list of [torch generator(s)](https://pytorch.org/docs/stable/generated/torch.Generator.html) | |
to make generation deterministic. | |
latents (`torch.FloatTensor`, *optional*): | |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image | |
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents | |
tensor will ge generated by sampling using the supplied random `generator`. | |
prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not | |
provided, text embeddings will be generated from `prompt` input argument. | |
negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input | |
argument. | |
output_type (`str`, *optional*, defaults to `"pil"`): | |
The output format of the generate image. Choose between | |
[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`. | |
return_dict (`bool`, *optional*, defaults to `True`): | |
Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] instead of a | |
plain tuple. | |
callback (`Callable`, *optional*): | |
A function that will be called every `callback_steps` steps during inference. The function will be | |
called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`. | |
callback_steps (`int`, *optional*, defaults to 1): | |
The frequency at which the `callback` function will be called. If not specified, the callback will be | |
called at every step. | |
cross_attention_kwargs (`dict`, *optional*): | |
A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under | |
`self.processor` in | |
[diffusers.models.attention_processor](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py). | |
controlnet_conditioning_scale (`float` or `List[float]`, *optional*, defaults to 0.5): | |
The outputs of the controlnet are multiplied by `controlnet_conditioning_scale` before they are added | |
to the residual in the original unet. If multiple ControlNets are specified in init, you can set the | |
corresponding scale as a list. Note that by default, we use a smaller conditioning scale for inpainting | |
than for [`~StableDiffusionControlNetPipeline.__call__`]. | |
guess_mode (`bool`, *optional*, defaults to `False`): | |
In this mode, the ControlNet encoder will try best to recognize the content of the input image even if | |
you remove all prompts. The `guidance_scale` between 3.0 and 5.0 is recommended. | |
control_guidance_start (`float` or `List[float]`, *optional*, defaults to 0.0): | |
The percentage of total steps at which the controlnet starts applying. | |
control_guidance_end (`float` or `List[float]`, *optional*, defaults to 1.0): | |
The percentage of total steps at which the controlnet stops applying. | |
Examples: | |
Returns: | |
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] or `tuple`: | |
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] if `return_dict` is True, otherwise a `tuple. | |
When returning a tuple, the first element is a list with the generated images, and the second element is a | |
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work" | |
(nsfw) content, according to the `safety_checker`. | |
""" | |
controlnet = self.controlnet._orig_mod if is_compiled_module(self.controlnet) else self.controlnet | |
# 0. Default height and width to unet | |
height, width = self._default_height_width(height, width, image) | |
prompt = promptA | |
negative_prompt = negative_promptA | |
# align format for control guidance | |
if not isinstance(control_guidance_start, list) and isinstance(control_guidance_end, list): | |
control_guidance_start = len(control_guidance_end) * [control_guidance_start] | |
elif not isinstance(control_guidance_end, list) and isinstance(control_guidance_start, list): | |
control_guidance_end = len(control_guidance_start) * [control_guidance_end] | |
elif not isinstance(control_guidance_start, list) and not isinstance(control_guidance_end, list): | |
mult = len(controlnet.nets) if isinstance(controlnet, MultiControlNetModel) else 1 | |
control_guidance_start, control_guidance_end = mult * [control_guidance_start], mult * [ | |
control_guidance_end | |
] | |
# 1. Check inputs. Raise error if not correct | |
self.check_inputs( | |
prompt, | |
control_image, | |
height, | |
width, | |
callback_steps, | |
negative_prompt, | |
prompt_embeds, | |
negative_prompt_embeds, | |
controlnet_conditioning_scale, | |
control_guidance_start, | |
control_guidance_end, | |
) | |
# 2. Define call parameters | |
if prompt is not None and isinstance(prompt, str): | |
batch_size = 1 | |
elif prompt is not None and isinstance(prompt, list): | |
batch_size = len(prompt) | |
else: | |
batch_size = prompt_embeds.shape[0] | |
device = self._execution_device | |
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2) | |
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1` | |
# corresponds to doing no classifier free guidance. | |
do_classifier_free_guidance = guidance_scale > 1.0 | |
if isinstance(controlnet, MultiControlNetModel) and isinstance(controlnet_conditioning_scale, float): | |
controlnet_conditioning_scale = [controlnet_conditioning_scale] * len(controlnet.nets) | |
global_pool_conditions = ( | |
controlnet.config.global_pool_conditions | |
if isinstance(controlnet, ControlNetModel) | |
else controlnet.nets[0].config.global_pool_conditions | |
) | |
guess_mode = guess_mode or global_pool_conditions | |
# 3. Encode input prompt | |
text_encoder_lora_scale = ( | |
cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None | |
) | |
prompt_embeds = self._encode_prompt( | |
promptA, | |
promptB, | |
tradoff, | |
device, | |
num_images_per_prompt, | |
do_classifier_free_guidance, | |
negative_promptA, | |
negative_promptB, | |
tradoff_nag, | |
prompt_embeds=prompt_embeds, | |
negative_prompt_embeds=negative_prompt_embeds, | |
lora_scale=text_encoder_lora_scale, | |
) | |
# 4. Prepare image | |
if isinstance(controlnet, ControlNetModel): | |
control_image = self.prepare_control_image( | |
image=control_image, | |
width=width, | |
height=height, | |
batch_size=batch_size * num_images_per_prompt, | |
num_images_per_prompt=num_images_per_prompt, | |
device=device, | |
dtype=controlnet.dtype, | |
do_classifier_free_guidance=do_classifier_free_guidance, | |
guess_mode=guess_mode, | |
) | |
elif isinstance(controlnet, MultiControlNetModel): | |
control_images = [] | |
for control_image_ in control_image: | |
control_image_ = self.prepare_control_image( | |
image=control_image_, | |
width=width, | |
height=height, | |
batch_size=batch_size * num_images_per_prompt, | |
num_images_per_prompt=num_images_per_prompt, | |
device=device, | |
dtype=controlnet.dtype, | |
do_classifier_free_guidance=do_classifier_free_guidance, | |
guess_mode=guess_mode, | |
) | |
control_images.append(control_image_) | |
control_image = control_images | |
else: | |
assert False | |
# 4. Preprocess mask and image - resizes image and mask w.r.t height and width | |
mask, masked_image, init_image = prepare_mask_and_masked_image( | |
image, mask_image, height, width, return_image=True | |
) | |
# 5. Prepare timesteps | |
self.scheduler.set_timesteps(num_inference_steps, device=device) | |
timesteps, num_inference_steps = self.get_timesteps( | |
num_inference_steps=num_inference_steps, strength=strength, device=device | |
) | |
# at which timestep to set the initial noise (n.b. 50% if strength is 0.5) | |
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt) | |
# create a boolean to check if the strength is set to 1. if so then initialise the latents with pure noise | |
is_strength_max = strength == 1.0 | |
# 6. Prepare latent variables | |
num_channels_latents = self.vae.config.latent_channels | |
num_channels_unet = self.unet.config.in_channels | |
return_image_latents = num_channels_unet == 4 | |
latents_outputs = self.prepare_latents( | |
batch_size * num_images_per_prompt, | |
num_channels_latents, | |
height, | |
width, | |
prompt_embeds.dtype, | |
device, | |
generator, | |
latents, | |
image=init_image, | |
timestep=latent_timestep, | |
is_strength_max=is_strength_max, | |
return_noise=True, | |
return_image_latents=return_image_latents, | |
) | |
if return_image_latents: | |
latents, noise, image_latents = latents_outputs | |
else: | |
latents, noise = latents_outputs | |
# 7. Prepare mask latent variables | |
mask, masked_image_latents = self.prepare_mask_latents( | |
mask, | |
masked_image, | |
batch_size * num_images_per_prompt, | |
height, | |
width, | |
prompt_embeds.dtype, | |
device, | |
generator, | |
do_classifier_free_guidance, | |
) | |
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline | |
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta) | |
# 7.1 Create tensor stating which controlnets to keep | |
controlnet_keep = [] | |
for i in range(len(timesteps)): | |
keeps = [ | |
1.0 - float(i / len(timesteps) < s or (i + 1) / len(timesteps) > e) | |
for s, e in zip(control_guidance_start, control_guidance_end) | |
] | |
controlnet_keep.append(keeps[0] if isinstance(controlnet, ControlNetModel) else keeps) | |
# 8. Denoising loop | |
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order | |
with self.progress_bar(total=num_inference_steps) as progress_bar: | |
for i, t in enumerate(timesteps): | |
# expand the latents if we are doing classifier free guidance | |
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents | |
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t) | |
# controlnet(s) inference | |
if guess_mode and do_classifier_free_guidance: | |
# Infer ControlNet only for the conditional batch. | |
control_model_input = latents | |
control_model_input = self.scheduler.scale_model_input(control_model_input, t) | |
controlnet_prompt_embeds = prompt_embeds.chunk(2)[1] | |
else: | |
control_model_input = latent_model_input | |
controlnet_prompt_embeds = prompt_embeds | |
if isinstance(controlnet_keep[i], list): | |
cond_scale = [c * s for c, s in zip(controlnet_conditioning_scale, controlnet_keep[i])] | |
else: | |
controlnet_cond_scale = controlnet_conditioning_scale | |
if isinstance(controlnet_cond_scale, list): | |
controlnet_cond_scale = controlnet_cond_scale[0] | |
cond_scale = controlnet_cond_scale * controlnet_keep[i] | |
down_block_res_samples, mid_block_res_sample = self.controlnet( | |
control_model_input, | |
t, | |
encoder_hidden_states=controlnet_prompt_embeds, | |
controlnet_cond=control_image, | |
conditioning_scale=cond_scale, | |
guess_mode=guess_mode, | |
return_dict=False, | |
) | |
if guess_mode and do_classifier_free_guidance: | |
# Infered ControlNet only for the conditional batch. | |
# To apply the output of ControlNet to both the unconditional and conditional batches, | |
# add 0 to the unconditional batch to keep it unchanged. | |
down_block_res_samples = [torch.cat([torch.zeros_like(d), d]) for d in down_block_res_samples] | |
mid_block_res_sample = torch.cat([torch.zeros_like(mid_block_res_sample), mid_block_res_sample]) | |
# predict the noise residual | |
if num_channels_unet == 9: | |
latent_model_input = torch.cat([latent_model_input, mask, masked_image_latents], dim=1) | |
noise_pred = self.unet( | |
latent_model_input, | |
t, | |
encoder_hidden_states=prompt_embeds, | |
cross_attention_kwargs=cross_attention_kwargs, | |
down_block_additional_residuals=down_block_res_samples, | |
mid_block_additional_residual=mid_block_res_sample, | |
return_dict=False, | |
)[0] | |
# perform guidance | |
if do_classifier_free_guidance: | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
# compute the previous noisy sample x_t -> x_t-1 | |
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0] | |
if num_channels_unet == 4: | |
init_latents_proper = image_latents[:1] | |
init_mask = mask[:1] | |
if i < len(timesteps) - 1: | |
noise_timestep = timesteps[i + 1] | |
init_latents_proper = self.scheduler.add_noise( | |
init_latents_proper, noise, torch.tensor([noise_timestep]) | |
) | |
latents = (1 - init_mask) * init_latents_proper + init_mask * latents | |
# call the callback, if provided | |
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0): | |
progress_bar.update() | |
if callback is not None and i % callback_steps == 0: | |
callback(i, t, latents) | |
# If we do sequential model offloading, let's offload unet and controlnet | |
# manually for max memory savings | |
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None: | |
self.unet.to("cpu") | |
self.controlnet.to("cpu") | |
torch.cuda.empty_cache() | |
if not output_type == "latent": | |
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0] | |
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype) | |
else: | |
image = latents | |
has_nsfw_concept = None | |
if has_nsfw_concept is None: | |
do_denormalize = [True] * image.shape[0] | |
else: | |
do_denormalize = [not has_nsfw for has_nsfw in has_nsfw_concept] | |
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize) | |
# Offload last model to CPU | |
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None: | |
self.final_offload_hook.offload() | |
if not return_dict: | |
return (image, has_nsfw_concept) | |
return StableDiffusionPipelineOutput(images=image, nsfw_content_detected=has_nsfw_concept) | |