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LICENSE ADDED
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+ Copyright (c) 2022 Stability AI and contributors
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+ CreativeML Open RAIL++-M License
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+ dated November 24, 2022
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+ Section I: PREAMBLE
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+ Multimodal generative models are being widely adopted and used, and have the potential to transform the way artists, among other individuals, conceive and benefit from AI or ML technologies as a tool for content creation.
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+ END OF TERMS AND CONDITIONS
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+ Attachment A
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+
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+ Use Restrictions
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+
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+ You agree not to use the Model or Derivatives of the Model:
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+
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+ - In any way that violates any applicable national, federal, state, local or international law or regulation;
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+ - For the purpose of exploiting, harming or attempting to exploit or harm minors in any way;
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+ - To generate or disseminate verifiably false information and/or content with the purpose of harming others;
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+ - To generate or disseminate personal identifiable information that can be used to harm an individual;
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+ - To defame, disparage or otherwise harass others;
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+ - For fully automated decision making that adversely impacts an individual’s legal rights or otherwise creates or modifies a binding, enforceable obligation;
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+ - For any use intended to or which has the effect of discriminating against or harming individuals or groups based on online or offline social behavior or known or predicted personal or personality characteristics;
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+ - To exploit any of the vulnerabilities of a specific group of persons based on their age, social, physical or mental characteristics, in order to materially distort the behavior of a person pertaining to that group in a manner that causes or is likely to cause that person or another person physical or psychological harm;
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+ - For any use intended to or which has the effect of discriminating against individuals or groups based on legally protected characteristics or categories;
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+ - To provide medical advice and medical results interpretation;
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+ - To generate or disseminate information for the purpose to be used for administration of justice, law enforcement, immigration or asylum processes, such as predicting an individual will commit fraud/crime commitment (e.g. by text profiling, drawing causal relationships between assertions made in documents, indiscriminate and arbitrarily-targeted use).
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README.md CHANGED
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- ---
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- title: Stable Difusion Scunge V1
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- emoji: 🚀
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- colorFrom: gray
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- colorTo: purple
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- sdk: gradio
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- sdk_version: 4.7.1
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- app_file: app.py
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- pinned: false
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- ---
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-
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- Check out the configuration reference at https://huggingface.co/docs/hub/spaces-config-reference
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
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+ # Stable Diffusion Version 2
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+ ![t2i](assets/stable-samples/txt2img/768/merged-0006.png)
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+ ![t2i](assets/stable-samples/txt2img/768/merged-0002.png)
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+ ![t2i](assets/stable-samples/txt2img/768/merged-0005.png)
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+
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+ This repository contains [Stable Diffusion](https://github.com/CompVis/stable-diffusion) models trained from scratch and will be continuously updated with
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+ new checkpoints. The following list provides an overview of all currently available models. More coming soon.
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+
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+ ## News
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+
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+ **March 24, 2023**
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+
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+ *Stable UnCLIP 2.1*
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+
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+ - New stable diffusion finetune (_Stable unCLIP 2.1_, [Hugging Face](https://huggingface.co/stabilityai/)) at 768x768 resolution, based on SD2.1-768. This model allows for image variations and mixing operations as described in [*Hierarchical Text-Conditional Image Generation with CLIP Latents*](https://arxiv.org/abs/2204.06125), and, thanks to its modularity, can be combined with other models such as [KARLO](https://github.com/kakaobrain/karlo). Comes in two variants: [*Stable unCLIP-L*](https://huggingface.co/stabilityai/stable-diffusion-2-1-unclip/blob/main/sd21-unclip-l.ckpt) and [*Stable unCLIP-H*](https://huggingface.co/stabilityai/stable-diffusion-2-1-unclip/blob/main/sd21-unclip-h.ckpt), which are conditioned on CLIP ViT-L and ViT-H image embeddings, respectively. Instructions are available [here](doc/UNCLIP.MD).
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+
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+ - A public demo of SD-unCLIP is already available at [clipdrop.co/stable-diffusion-reimagine](https://clipdrop.co/stable-diffusion-reimagine)
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+
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+
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+ **December 7, 2022**
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+
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+ *Version 2.1*
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+
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+ - New stable diffusion model (_Stable Diffusion 2.1-v_, [Hugging Face](https://huggingface.co/stabilityai/stable-diffusion-2-1)) at 768x768 resolution and (_Stable Diffusion 2.1-base_, [HuggingFace](https://huggingface.co/stabilityai/stable-diffusion-2-1-base)) at 512x512 resolution, both based on the same number of parameters and architecture as 2.0 and fine-tuned on 2.0, on a less restrictive NSFW filtering of the [LAION-5B](https://laion.ai/blog/laion-5b/) dataset.
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+ Per default, the attention operation of the model is evaluated at full precision when `xformers` is not installed. To enable fp16 (which can cause numerical instabilities with the vanilla attention module on the v2.1 model) , run your script with `ATTN_PRECISION=fp16 python <thescript.py>`
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+
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+ **November 24, 2022**
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+
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+ *Version 2.0*
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+
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+ - New stable diffusion model (_Stable Diffusion 2.0-v_) at 768x768 resolution. Same number of parameters in the U-Net as 1.5, but uses [OpenCLIP-ViT/H](https://github.com/mlfoundations/open_clip) as the text encoder and is trained from scratch. _SD 2.0-v_ is a so-called [v-prediction](https://arxiv.org/abs/2202.00512) model.
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+ - The above model is finetuned from _SD 2.0-base_, which was trained as a standard noise-prediction model on 512x512 images and is also made available.
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+ - Added a [x4 upscaling latent text-guided diffusion model](#image-upscaling-with-stable-diffusion).
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+ - New [depth-guided stable diffusion model](#depth-conditional-stable-diffusion), finetuned from _SD 2.0-base_. The model is conditioned on monocular depth estimates inferred via [MiDaS](https://github.com/isl-org/MiDaS) and can be used for structure-preserving img2img and shape-conditional synthesis.
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+
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+ ![d2i](assets/stable-samples/depth2img/depth2img01.png)
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+ - A [text-guided inpainting model](#image-inpainting-with-stable-diffusion), finetuned from SD _2.0-base_.
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+
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+ We follow the [original repository](https://github.com/CompVis/stable-diffusion) and provide basic inference scripts to sample from the models.
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+
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+ ________________
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+ *The original Stable Diffusion model was created in a collaboration with [CompVis](https://arxiv.org/abs/2202.00512) and [RunwayML](https://runwayml.com/) and builds upon the work:*
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+
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+ [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://ommer-lab.com/research/latent-diffusion-models/)<br/>
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+ [Robin Rombach](https://github.com/rromb)\*,
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+ [Andreas Blattmann](https://github.com/ablattmann)\*,
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+ [Dominik Lorenz](https://github.com/qp-qp)\,
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+ [Patrick Esser](https://github.com/pesser),
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+ [Björn Ommer](https://hci.iwr.uni-heidelberg.de/Staff/bommer)<br/>
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+ _[CVPR '22 Oral](https://openaccess.thecvf.com/content/CVPR2022/html/Rombach_High-Resolution_Image_Synthesis_With_Latent_Diffusion_Models_CVPR_2022_paper.html) |
52
+ [GitHub](https://github.com/CompVis/latent-diffusion) | [arXiv](https://arxiv.org/abs/2112.10752) | [Project page](https://ommer-lab.com/research/latent-diffusion-models/)_
53
+
54
+ and [many others](#shout-outs).
55
+
56
+ Stable Diffusion is a latent text-to-image diffusion model.
57
+ ________________________________
58
+
59
+ ## Requirements
60
+
61
+ You can update an existing [latent diffusion](https://github.com/CompVis/latent-diffusion) environment by running
62
+
63
+ ```
64
+ conda install pytorch==1.12.1 torchvision==0.13.1 -c pytorch
65
+ pip install transformers==4.19.2 diffusers invisible-watermark
66
+ pip install -e .
67
+ ```
68
+ #### xformers efficient attention
69
+ For more efficiency and speed on GPUs,
70
+ we highly recommended installing the [xformers](https://github.com/facebookresearch/xformers)
71
+ library.
72
+
73
+ Tested on A100 with CUDA 11.4.
74
+ Installation needs a somewhat recent version of nvcc and gcc/g++, obtain those, e.g., via
75
+ ```commandline
76
+ export CUDA_HOME=/usr/local/cuda-11.4
77
+ conda install -c nvidia/label/cuda-11.4.0 cuda-nvcc
78
+ conda install -c conda-forge gcc
79
+ conda install -c conda-forge gxx_linux-64==9.5.0
80
+ ```
81
+
82
+ Then, run the following (compiling takes up to 30 min).
83
+
84
+ ```commandline
85
+ cd ..
86
+ git clone https://github.com/facebookresearch/xformers.git
87
+ cd xformers
88
+ git submodule update --init --recursive
89
+ pip install -r requirements.txt
90
+ pip install -e .
91
+ cd ../stablediffusion
92
+ ```
93
+ Upon successful installation, the code will automatically default to [memory efficient attention](https://github.com/facebookresearch/xformers)
94
+ for the self- and cross-attention layers in the U-Net and autoencoder.
95
+
96
+ ## General Disclaimer
97
+ Stable Diffusion models are general text-to-image diffusion models and therefore mirror biases and (mis-)conceptions that are present
98
+ in their training data. Although efforts were made to reduce the inclusion of explicit pornographic material, **we do not recommend using the provided weights for services or products without additional safety mechanisms and considerations.
99
+ The weights are research artifacts and should be treated as such.**
100
+ Details on the training procedure and data, as well as the intended use of the model can be found in the corresponding [model card](https://huggingface.co/stabilityai/stable-diffusion-2).
101
+ The weights are available via [the StabilityAI organization at Hugging Face](https://huggingface.co/StabilityAI) under the [CreativeML Open RAIL++-M License](LICENSE-MODEL).
102
+
103
+
104
+
105
+ ## Stable Diffusion v2
106
+
107
+ Stable Diffusion v2 refers to a specific configuration of the model
108
+ architecture that uses a downsampling-factor 8 autoencoder with an 865M UNet
109
+ and OpenCLIP ViT-H/14 text encoder for the diffusion model. The _SD 2-v_ model produces 768x768 px outputs.
110
+
111
+ Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
112
+ 5.0, 6.0, 7.0, 8.0) and 50 DDIM sampling steps show the relative improvements of the checkpoints:
113
+
114
+ ![sd evaluation results](assets/model-variants.jpg)
115
+
116
+
117
+
118
+ ### Text-to-Image
119
+ ![txt2img-stable2](assets/stable-samples/txt2img/merged-0003.png)
120
+ ![txt2img-stable2](assets/stable-samples/txt2img/merged-0001.png)
121
+
122
+ Stable Diffusion 2 is a latent diffusion model conditioned on the penultimate text embeddings of a CLIP ViT-H/14 text encoder.
123
+ We provide a [reference script for sampling](#reference-sampling-script).
124
+ #### Reference Sampling Script
125
+
126
+ This script incorporates an [invisible watermarking](https://github.com/ShieldMnt/invisible-watermark) of the outputs, to help viewers [identify the images as machine-generated](scripts/tests/test_watermark.py).
127
+ We provide the configs for the _SD2-v_ (768px) and _SD2-base_ (512px) model.
128
+
129
+ First, download the weights for [_SD2.1-v_](https://huggingface.co/stabilityai/stable-diffusion-2-1) and [_SD2.1-base_](https://huggingface.co/stabilityai/stable-diffusion-2-1-base).
130
+
131
+ To sample from the _SD2.1-v_ model, run the following:
132
+
133
+ ```
134
+ python scripts/txt2img.py --prompt "a professional photograph of an astronaut riding a horse" --ckpt <path/to/768model.ckpt/> --config configs/stable-diffusion/v2-inference-v.yaml --H 768 --W 768
135
+ ```
136
+ or try out the Web Demo: [![Hugging Face Spaces](https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue)](https://huggingface.co/spaces/stabilityai/stable-diffusion).
137
+
138
+ To sample from the base model, use
139
+ ```
140
+ python scripts/txt2img.py --prompt "a professional photograph of an astronaut riding a horse" --ckpt <path/to/model.ckpt/> --config <path/to/config.yaml/>
141
+ ```
142
+
143
+ By default, this uses the [DDIM sampler](https://arxiv.org/abs/2010.02502), and renders images of size 768x768 (which it was trained on) in 50 steps.
144
+ Empirically, the v-models can be sampled with higher guidance scales.
145
+
146
+ Note: The inference config for all model versions is designed to be used with EMA-only checkpoints.
147
+ For this reason `use_ema=False` is set in the configuration, otherwise the code will try to switch from
148
+ non-EMA to EMA weights.
149
+
150
+ #### Enable Intel® Extension for PyTorch* optimizations in Text-to-Image script
151
+
152
+ If you're planning on running Text-to-Image on Intel® CPU, try to sample an image with TorchScript and Intel® Extension for PyTorch* optimizations. Intel® Extension for PyTorch* extends PyTorch by enabling up-to-date features optimizations for an extra performance boost on Intel® hardware. It can optimize memory layout of the operators to Channel Last memory format, which is generally beneficial for Intel CPUs, take advantage of the most advanced instruction set available on a machine, optimize operators and many more.
153
+
154
+ **Prerequisites**
155
+
156
+ Before running the script, make sure you have all needed libraries installed. (the optimization was checked on `Ubuntu 20.04`). Install [jemalloc](https://github.com/jemalloc/jemalloc), [numactl](https://linux.die.net/man/8/numactl), Intel® OpenMP and Intel® Extension for PyTorch*.
157
+
158
+ ```bash
159
+ apt-get install numactl libjemalloc-dev
160
+ pip install intel-openmp
161
+ pip install intel_extension_for_pytorch -f https://software.intel.com/ipex-whl-stable
162
+ ```
163
+
164
+ To sample from the _SD2.1-v_ model with TorchScript+IPEX optimizations, run the following. Remember to specify desired number of instances you want to run the program on ([more](https://github.com/intel/intel-extension-for-pytorch/blob/master/intel_extension_for_pytorch/cpu/launch.py#L48)).
165
+
166
+ ```
167
+ MALLOC_CONF=oversize_threshold:1,background_thread:true,metadata_thp:auto,dirty_decay_ms:9000000000,muzzy_decay_ms:9000000000 python -m intel_extension_for_pytorch.cpu.launch --ninstance <number of an instance> --enable_jemalloc scripts/txt2img.py --prompt \"a corgi is playing guitar, oil on canvas\" --ckpt <path/to/768model.ckpt/> --config configs/stable-diffusion/intel/v2-inference-v-fp32.yaml --H 768 --W 768 --precision full --device cpu --torchscript --ipex
168
+ ```
169
+
170
+ To sample from the base model with IPEX optimizations, use
171
+
172
+ ```
173
+ MALLOC_CONF=oversize_threshold:1,background_thread:true,metadata_thp:auto,dirty_decay_ms:9000000000,muzzy_decay_ms:9000000000 python -m intel_extension_for_pytorch.cpu.launch --ninstance <number of an instance> --enable_jemalloc scripts/txt2img.py --prompt \"a corgi is playing guitar, oil on canvas\" --ckpt <path/to/model.ckpt/> --config configs/stable-diffusion/intel/v2-inference-fp32.yaml --n_samples 1 --n_iter 4 --precision full --device cpu --torchscript --ipex
174
+ ```
175
+
176
+ If you're using a CPU that supports `bfloat16`, consider sample from the model with bfloat16 enabled for a performance boost, like so
177
+
178
+ ```bash
179
+ # SD2.1-v
180
+ MALLOC_CONF=oversize_threshold:1,background_thread:true,metadata_thp:auto,dirty_decay_ms:9000000000,muzzy_decay_ms:9000000000 python -m intel_extension_for_pytorch.cpu.launch --ninstance <number of an instance> --enable_jemalloc scripts/txt2img.py --prompt \"a corgi is playing guitar, oil on canvas\" --ckpt <path/to/768model.ckpt/> --config configs/stable-diffusion/intel/v2-inference-v-bf16.yaml --H 768 --W 768 --precision full --device cpu --torchscript --ipex --bf16
181
+ # SD2.1-base
182
+ MALLOC_CONF=oversize_threshold:1,background_thread:true,metadata_thp:auto,dirty_decay_ms:9000000000,muzzy_decay_ms:9000000000 python -m intel_extension_for_pytorch.cpu.launch --ninstance <number of an instance> --enable_jemalloc scripts/txt2img.py --prompt \"a corgi is playing guitar, oil on canvas\" --ckpt <path/to/model.ckpt/> --config configs/stable-diffusion/intel/v2-inference-bf16.yaml --precision full --device cpu --torchscript --ipex --bf16
183
+ ```
184
+
185
+ ### Image Modification with Stable Diffusion
186
+
187
+ ![depth2img-stable2](assets/stable-samples/depth2img/merged-0000.png)
188
+ #### Depth-Conditional Stable Diffusion
189
+
190
+ To augment the well-established [img2img](https://github.com/CompVis/stable-diffusion#image-modification-with-stable-diffusion) functionality of Stable Diffusion, we provide a _shape-preserving_ stable diffusion model.
191
+
192
+
193
+ Note that the original method for image modification introduces significant semantic changes w.r.t. the initial image.
194
+ If that is not desired, download our [depth-conditional stable diffusion](https://huggingface.co/stabilityai/stable-diffusion-2-depth) model and the `dpt_hybrid` MiDaS [model weights](https://github.com/intel-isl/DPT/releases/download/1_0/dpt_hybrid-midas-501f0c75.pt), place the latter in a folder `midas_models` and sample via
195
+ ```
196
+ python scripts/gradio/depth2img.py configs/stable-diffusion/v2-midas-inference.yaml <path-to-ckpt>
197
+ ```
198
+
199
+ or
200
+
201
+ ```
202
+ streamlit run scripts/streamlit/depth2img.py configs/stable-diffusion/v2-midas-inference.yaml <path-to-ckpt>
203
+ ```
204
+
205
+ This method can be used on the samples of the base model itself.
206
+ For example, take [this sample](assets/stable-samples/depth2img/old_man.png) generated by an anonymous discord user.
207
+ Using the [gradio](https://gradio.app) or [streamlit](https://streamlit.io/) script `depth2img.py`, the MiDaS model first infers a monocular depth estimate given this input,
208
+ and the diffusion model is then conditioned on the (relative) depth output.
209
+
210
+ <p align="center">
211
+ <b> depth2image </b><br/>
212
+ <img src=assets/stable-samples/depth2img/d2i.gif>
213
+ </p>
214
+
215
+ This model is particularly useful for a photorealistic style; see the [examples](assets/stable-samples/depth2img).
216
+ For a maximum strength of 1.0, the model removes all pixel-based information and only relies on the text prompt and the inferred monocular depth estimate.
217
+
218
+ ![depth2img-stable3](assets/stable-samples/depth2img/merged-0005.png)
219
+
220
+ #### Classic Img2Img
221
+
222
+ For running the "classic" img2img, use
223
+ ```
224
+ python scripts/img2img.py --prompt "A fantasy landscape, trending on artstation" --init-img <path-to-img.jpg> --strength 0.8 --ckpt <path/to/model.ckpt>
225
+ ```
226
+ and adapt the checkpoint and config paths accordingly.
227
+
228
+ ### Image Upscaling with Stable Diffusion
229
+ ![upscaling-x4](assets/stable-samples/upscaling/merged-dog.png)
230
+ After [downloading the weights](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler), run
231
+ ```
232
+ python scripts/gradio/superresolution.py configs/stable-diffusion/x4-upscaling.yaml <path-to-checkpoint>
233
+ ```
234
+
235
+ or
236
+
237
+ ```
238
+ streamlit run scripts/streamlit/superresolution.py -- configs/stable-diffusion/x4-upscaling.yaml <path-to-checkpoint>
239
+ ```
240
+
241
+ for a Gradio or Streamlit demo of the text-guided x4 superresolution model.
242
+ This model can be used both on real inputs and on synthesized examples. For the latter, we recommend setting a higher
243
+ `noise_level`, e.g. `noise_level=100`.
244
+
245
+ ### Image Inpainting with Stable Diffusion
246
+
247
+ ![inpainting-stable2](assets/stable-inpainting/merged-leopards.png)
248
+
249
+ [Download the SD 2.0-inpainting checkpoint](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) and run
250
+
251
+ ```
252
+ python scripts/gradio/inpainting.py configs/stable-diffusion/v2-inpainting-inference.yaml <path-to-checkpoint>
253
+ ```
254
+
255
+ or
256
+
257
+ ```
258
+ streamlit run scripts/streamlit/inpainting.py -- configs/stable-diffusion/v2-inpainting-inference.yaml <path-to-checkpoint>
259
+ ```
260
+
261
+ for a Gradio or Streamlit demo of the inpainting model.
262
+ This scripts adds invisible watermarking to the demo in the [RunwayML](https://github.com/runwayml/stable-diffusion/blob/main/scripts/inpaint_st.py) repository, but both should work interchangeably with the checkpoints/configs.
263
+
264
+
265
+
266
+ ## Shout-Outs
267
+ - Thanks to [Hugging Face](https://huggingface.co/) and in particular [Apolinário](https://github.com/apolinario) for support with our model releases!
268
+ - Stable Diffusion would not be possible without [LAION](https://laion.ai/) and their efforts to create open, large-scale datasets.
269
+ - The [DeepFloyd team](https://twitter.com/deepfloydai) at Stability AI, for creating the subset of [LAION-5B](https://laion.ai/blog/laion-5b/) dataset used to train the model.
270
+ - Stable Diffusion 2.0 uses [OpenCLIP](https://laion.ai/blog/large-openclip/), trained by [Romain Beaumont](https://github.com/rom1504).
271
+ - Our codebase for the diffusion models builds heavily on [OpenAI's ADM codebase](https://github.com/openai/guided-diffusion)
272
+ and [https://github.com/lucidrains/denoising-diffusion-pytorch](https://github.com/lucidrains/denoising-diffusion-pytorch).
273
+ Thanks for open-sourcing!
274
+ - [CompVis](https://github.com/CompVis/stable-diffusion) initial stable diffusion release
275
+ - [Patrick](https://github.com/pesser)'s [implementation](https://github.com/runwayml/stable-diffusion/blob/main/scripts/inpaint_st.py) of the streamlit demo for inpainting.
276
+ - `img2img` is an application of [SDEdit](https://arxiv.org/abs/2108.01073) by [Chenlin Meng](https://cs.stanford.edu/~chenlin/) from the [Stanford AI Lab](https://cs.stanford.edu/~ermon/website/).
277
+ - [Kat's implementation]((https://github.com/CompVis/latent-diffusion/pull/51)) of the [PLMS](https://arxiv.org/abs/2202.09778) sampler, and [more](https://github.com/crowsonkb/k-diffusion).
278
+ - [DPMSolver](https://arxiv.org/abs/2206.00927) [integration](https://github.com/CompVis/stable-diffusion/pull/440) by [Cheng Lu](https://github.com/LuChengTHU).
279
+ - Facebook's [xformers](https://github.com/facebookresearch/xformers) for efficient attention computation.
280
+ - [MiDaS](https://github.com/isl-org/MiDaS) for monocular depth estimation.
281
+
282
+
283
+ ## License
284
+
285
+ The code in this repository is released under the MIT License.
286
+
287
+ The weights are available via [the StabilityAI organization at Hugging Face](https://huggingface.co/StabilityAI), and released under the [CreativeML Open RAIL++-M License](LICENSE-MODEL) License.
288
+
289
+ ## BibTeX
290
+
291
+ ```
292
+ @misc{rombach2021highresolution,
293
+ title={High-Resolution Image Synthesis with Latent Diffusion Models},
294
+ author={Robin Rombach and Andreas Blattmann and Dominik Lorenz and Patrick Esser and Björn Ommer},
295
+ year={2021},
296
+ eprint={2112.10752},
297
+ archivePrefix={arXiv},
298
+ primaryClass={cs.CV}
299
+ }
300
+ ```
301
+
302
+
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checkpoints/checkpoints.txt ADDED
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+ Put unCLIP checkpoints here.
configs/karlo/decoder_900M_vit_l.yaml ADDED
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+ model:
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+ type: t2i-decoder
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+ diffusion_sampler: uniform
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+ hparams:
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+ image_size: 64
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+ num_channels: 320
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+ num_res_blocks: 3
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+ channel_mult: ''
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+ attention_resolutions: 32,16,8
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+ num_heads: -1
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+ num_head_channels: 64
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+ num_heads_upsample: -1
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+ use_scale_shift_norm: true
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+ dropout: 0.1
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+ clip_dim: 768
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+ clip_emb_mult: 4
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+ text_ctx: 77
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+ xf_width: 1536
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+ xf_layers: 0
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+ xf_heads: 0
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+ xf_final_ln: false
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+ resblock_updown: true
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+ learn_sigma: true
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+ text_drop: 0.3
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+ clip_emb_type: image
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+ clip_emb_drop: 0.1
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+ use_plm: true
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+
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+ diffusion:
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+ steps: 1000
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+ learn_sigma: true
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+ sigma_small: false
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+ noise_schedule: squaredcos_cap_v2
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+ use_kl: false
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+ predict_xstart: false
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+ rescale_learned_sigmas: true
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+ timestep_respacing: ''