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--- |
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license: other |
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pipeline_tag: image-to-image |
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--- |
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# StableSR Model Card |
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This model card focuses on the models associated with the StableSR, available [here](https://github.com/IceClear/StableSR). |
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## Model Details |
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- **Developed by:** Jianyi Wang |
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- **Model type:** Diffusion-based image super-resolution model |
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- **License:** [S-Lab License 1.0](https://github.com/IceClear/StableSR/blob/main/LICENSE.txt) |
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- **Model Description:** This is the model used in [Paper](https://arxiv.org/abs/2305.07015). |
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- **Resources for more information:** [GitHub Repository](https://github.com/IceClear/StableSR). |
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- **Cite as:** |
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@InProceedings{wang2023exploiting, |
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author = {Wang, Jianyi and Yue, Zongsheng and Zhou, Shangchen and Chan, Kelvin CK and Loy, Chen Change}, |
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title = {Exploiting Diffusion Prior for Real-World Image Super-Resolution}, |
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booktitle = {arXiv preprint arXiv:2305.07015}, |
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year = {2023}, |
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} |
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# Uses |
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Please refer to [S-Lab License 1.0](https://github.com/IceClear/StableSR/blob/main/LICENSE.txt) |
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## Limitations and Bias |
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### Limitations |
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- StableSR still requires multiple steps for generating an image, which is much slower than GAN-based approaches, especially for large images beyond 512 or 768. |
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- StableSR sometimes cannot keep 100% fidelity due to its generative nature. |
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- StableSR sometimes cannot generate perfect details under complex real-world scenarios. |
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### Bias |
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While our model is based on a pre-trained Stable Diffusion model, currently we do not observe obvious bias in generated results. |
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We conjecture the main reason is that our model does not rely on text prompts but on low-resolution images. |
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Such strong conditions make our model less likely to be affected. |
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## Training |
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**Training Data** |
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The model developer used the following dataset for training the model: |
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- Our diffusion model is finetuned on DF2K (DIV2K and Flickr2K) + OST datasets, available [here](https://github.com/xinntao/Real-ESRGAN/blob/master/docs/Training.md). |
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- We further generate 100k synthetic LR-HR pairs on DF2K_OST using the finetuned diffusion model for training the CFW module. |
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**Training Procedure** |
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StableSR is an image super-resolution model finetuned on [Stable Diffusion](https://github.com/Stability-AI/stablediffusion), further equipped with a time-aware encoder and a controllable feature wrapping (CFW) module. |
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- Following Stable Diffusion, images are encoded through the fixed autoencoder, which turns images into latent representations. The autoencoder uses a relative downsampling factor of 8 and maps images of shape H x W x 3 to latents of shape H/f x W/f x 4. |
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- The latent representations are fed to the time-aware encoder as guidance. |
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- The loss is the same as Stable Diffusion. |
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- After finetuning the diffusion model, we further train the CFW module using the data generated by the finetuned diffusion model. |
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- The autoencoder model is fixed and only CFW is trainable. |
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- The loss is similar to training an autoencoder, except that we use a fixed adversarial loss weight of 0.025 rather than a self-adjustable one. |
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We currently provide the following checkpoints: |
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- [stablesr_000117.ckpt](https://huggingface.co/Iceclear/StableSR/resolve/main/stablesr_000117.ckpt): Diffusion model finetuned on [SD2.1-512base](https://huggingface.co/stabilityai/stable-diffusion-2-1-base) with DF2K_OST dataset for 117 epochs. |
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- [vqgan_cfw_00011.ckpt](https://huggingface.co/Iceclear/StableSR/resolve/main/vqgan_cfw_00011.ckpt): CFW module with fixed autoencoder trained on synthetic paired data for 11 epochs. |
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- [stablesr_768v_000139.ckpt](https://huggingface.co/Iceclear/StableSR/blob/main/stablesr_768v_000139.ckpt): Diffusion model finetuned on [SD2.1-768v](https://huggingface.co/stabilityai/stable-diffusion-2-1) with DF2K_OST dataset for 139 epochs. |
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## Evaluation Results |
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See [Paper](https://arxiv.org/abs/2305.07015) for details. |