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A typical adult does not brush his or her teeth adequately, nor does the typical adult floss his or her teeth with the frequency recommended by dentists. Gingivitis is usually caused by bacterial plaque that accumulates in the spaces between the gums and the teeth and in calculus that forms on the teeth. The bacteria produce foreign chemicals and toxins that cause inflammation of the gums around the teeth. This inflammation can cause deep pockets between the teeth and gums and loss of bone around teeth otherwise known as periodontitis. Since the bone in the jaws holds the teeth into the jaws, the loss of bone may cause teeth to become loose, and eventually fall out. If the inflammation in the gums becomes especially well-developed, it can invade the gums and allow tiny amounts of bacteria and bacterial toxins to enter the bloodstream. Studies suggest this can result in a generalized increase in inflammation in the body cause possible long term heart problems. Periodontitis has also been linked to diabetes, arterosclerosis, osteoporosis and pre-term low birth weight babies. Bacteria and bacterial toxins are carried by tooth brushes, and brushing alone may cause the germs to remain in the mouth if the tooth brush is not sanitized between uses. There is a need for a device that will hold and sanitize toothbrushes. The use of common antiseptic mouthwashes in combination with brushing, even if the brushing is otherwise inadequate, will kill harmful bacteria, and improve oral hygiene. The device should therefore encourage the use of mouthwash, and facilitate improved sanitation of toothbrushes.

The pest management practices throughout the world are increasingly leaning toward more environment friendly agriculture without disturbing the balance of the eco-system[@b1]. Thus, the evolution of modern organic agriculture depends on biological control agents to disrupt and control pest attacks which impart serious adverse effects on the agricultural economy[@b2]. Guavas (*Psidium guajava* Linn.) comprise plants that contain nearly hundred species of tropical shrubs and trees. Guavas are rich in dietary fibers, minerals, vitamins A and C and folic acid[@b3]. The fruit is native to central and northern South America, and cultivated throughout in South-east Asia, Florida, Hawaii, the Caribbean and Africa[@b3]. Fruitfly (*Bactrocera dorsalis* Hendel, [Figure 1a](#f1){ref-type="fig"}) is considered to be a harmful key pest of guava. It causes significant yield losses and quality degradations of the fruits particularly during the rainy season[@b4]. This results in a poor commercialization of the crops in domestic markets and imposes quarantine restrictions from the international markets[@b5]. *Bactrocera dorsalis* (*B. dorsalis*) are polyphagous[@b6] pests and it breeds on a number of other fruits and vegetables as well. For example, fruits such as, apple (*Annona squamosa*), banana (*Musa paradisiaca*), mango (*Mangifera indica*), orange (*Citrus sinensis*), peach (*Prunus persica*), plum (*Prunus domestica*) *etc*. are also affected by this fruit fly[@b7][@b8]. The fruit flies are attracted toward the concerned fruits and infect them ([Figure 1b](#f1){ref-type="fig"}). Attacked fruits show signs of ovipositional damage in the form of minute depressions which can be even seen from outside ([Figure 1c, d](#f1){ref-type="fig"}). The affected fruits soften at the infested positions, then disintegrate and eventually drop down prematurely ([Figure 1e](#f1){ref-type="fig"}). Eggs of *B. dorsalis* are laid below the skin of the host fruit and the maggots feed on the pulp of the fruits ([Figure 1f](#f1){ref-type="fig"}). These hatch within \~1--3 days and the larvae feed for another 9--35 days. The pupariation process occurs in the soil near the host plant and the adult flies emerge after \~1--2 weeks and persist throughout the year[@b9]. For the pest management of *B. dorsalis,* several methods have been developed in the last few decades. These pest management practices include destruction of the fallen and infested fruits[@b10], cold treatment[@b11], heat teatment[@b12][@b13], bait sprays method containing protein hydrolysate[@b14], biological control method[@b15], sterile insect techniques[@b16][@b17] *etc.* However, most of these treatments fail due to one or another type of drawbacks. Use of insecticides provides another alternative. But it affects both the agricultural fields, habitation and even the locality resulting in affecting non-target species such as humans, mammals, arthropods and other natural enemies[@b18]. In addition, the pest populations have increasingly become resistant to the chemical pesticides compelling the farmers to resort to an use of even larger amounts of pesticides[@b19][@b20][@b21]. These pose significant risk on the human health and the environment in a bid to achieve appropriate control of the pest. To alleviate such problems highly viscous hydrogels have been used to deliver pheromones[@b22]. However, hydrogels often swell and shrink with the changes in humidity and temperature. Therefore, such hydrogels cannot be used in all seasons and in particular during the rainy season. Male annihilation technique has been used in the pest management of *B. dorsalis* due to the attraction of males to methyl eugenol (ME) at very low concentrations[@b23][@b24]. However, there is a significant limitation in that these solutions contain pesticides and have a limited shelf-life. To overcome these problems a new yet simple approach is essential which should provide viable solutions in the pest management arena. Pheromones are naturally occurring volatile semiochemicals and are considered as eco-friendly biological control agents[@b25]. They induce impaired sexual communication and mating disruptions among pests[@b26]. High specificity of pheromones towards a given species further enables pest control without weakening the beneficial organisms[@b27]. However, due to the unstable nature of pheromones, such as photo-oxidation, auto-oxidation, isomerization and volatility *etc.*[@b28], the formulation development for the field trials involving pheromones needs a major refinement specifically in terms of the slow-release and protection from decomposition under ambient conditions[@b29]. ME in this context may be useful since it belongs to the class of a para-pheromone which mimics the role of a male attractant[@b30][@b31]. Using supramolecular self-assembly principles and nanotechnology we have been able to achieve an increased shelf-life of the pheromone (ME) when immobilized in a nanogel. We demonstrate here that such pheromone nanogels exhibit high residual activity, excellent efficacy in open orchard even during adverse seasons and are also environment friendly. Results ======= Nanotechnology plays a pivotal role in the increase of the effectiveness of pheromones in terms of their increased shelf-life through their entrapment by suitable host matrix[@b32]. Low-molecular mass gelators (LMMGs) have attracted much attention because of their interesting properties and potential uses in various fields of research and industry[@b33][@b34]. LMMGs such as all-*trans* tri(*p*-phenylenevinylene) bis-aldoxime (**1**)[@b35], self-assemble through weak intermolecular non-covalent interactions such as hydrogen bonding, π-π stacking and van der Waals leading to the gelation of specific solvents ([Figure 2](#f2){ref-type="fig"})[@b36][@b37][@b38][@b39]. In the present work, we demonstrate the nanogel formation upon entrapment and immobilization of ME into the entangled three-dimensional nano-sized supramolecular networks of **1**. Such nanogels provide high pheromone retention capacity, enhanced shelf-life and protection of ME from environmental decompositions (from exposure to ambient air, water and sunshine). ME alone however, is quite volatile on its own and evaporates under ambient conditions significantly rapidly during the summer (see below). Immobilization of methyl eugenol in nanogels -------------------------------------------- The present invention describes a simple and practical method by which pheromones could be immobilized in the nanogels. ME alone is a colorless semi-viscous liquid ([Figure 2b](#f2){ref-type="fig"}) which turns yellow ([Figure 2c](#f2){ref-type="fig"}) on solubilization of the gelator **1** (2.5 mg/mL) upon brief warming of the mixture to \~65°C. In this condition, the pheromone did not undergo any decomposition. When cooled to room temperature without any external perturbation for a period of 10--15 min, the yellow solution turned into a thick, viscous orange colored mass, which did not flow under the influence of gravity when the vial was inverted upside down demonstrating the onset of gelation ([Figure 2d](#f2){ref-type="fig"}). The gelation was retained unperturbed even after several consecutive heating and cooling cycles (from room temperature to 70°C). The nanogel was thermoreversible in nature, a property which was found to be useful for its long term use in the field. Thus the preparation of the nanogels (**1** + ME) is indeed convenient without requiring addition of any organic solvents or chemical cross-linking agents etc. Also, after several heat-cool cycles, or upon exposure to open orchard for a number of weeks, neither a thermal or photo-chemical decomposition of ME occurred as evidenced from GC-MS analysis (see, [Supplementary Information](#s1){ref-type="supplementary-material"}, \[[SI](#s1){ref-type="supplementary-material"}\] [Figure S1 and S2](#s1){ref-type="supplementary-material"}). This ensured good chemical stability of the resulting nanogels. The minimum amount of the gelator (**1**) required for the formation of the gel is known as the minimum gelator concentration (MGC). This was estimated to be \~2.5 mg/mL indicating high ME content in the nanogel. Thus as little as \~2.5 mg of **1** could immobilize about 1 mL of ME. The freeze-dried gel was further characterized using both transmission electron microscopy (TEM) and scanning electron microscopy (SEM). The presence of three-dimensional nano-sized fibrillar networks having fiber diameters ranging from 100--200 nm was evidenced. These fibers collate upon themselves in the nanogel to form aggregated bundles having the fiber diameters ranging from 500 nm to 1 μm ([Figure 3](#f3){ref-type="fig"}). Thermal stability of nanogels ----------------------------- The melting temperature of the nanogels, *i.e.* gel-to-sol transition (*T*~gel~) increased progressively with increasing concentration of **1**, a phenomenon that is general for many LMMGs ([Figure 4a](#f4){ref-type="fig"})[@b34][@b35]. This indicates that, with increasing concentration of the gelator, the density of the gel assembly increases ensuring participation of a larger number of gelator molecules per unit volume in the thermally induced transitions. At the concentration of 12 mg/mL of **1** in ME the *T*~gel~ reached \~63°C which was well above the ambient temperature even during peak summer in India. Thus the thermal stability of the sample should be adequate for the agricultural field trials even in hot climates around different geographical regions of the world. Rheological properties and mechanical stability of gels ------------------------------------------------------- The oscillatory frequency response of the nanogel revealed that the elastic modulus G′ is independent of frequency and much higher than the viscous modulus G″ over the entire angular frequency range (0.1--100 rad/s) at 0.01% strain amplitude ([Figure 4b](#f4){ref-type="fig"}). A \~13-fold higher G′ over G″ suggested a substantial elastic response of the gels that did not relax over the long timescale. This is a typical characteristic of the nanogels[@b34][@b35]. The viscoelastic solid-like nature of the gel increases with increasing concentration of **1** and it becomes sufficiently strong (G′ \~2800 Pa) at 12 mg/mL. Under oscillatory stress, the nanogels show G′ \> G″ by an order of magnitude indicating the dominant elastic behavior (stiffness) of the system and show the resistance of the mass to flow under applied stress ([Figure 4c](#f4){ref-type="fig"}). An applied stress above which the gel starts to flow is called the yield stress (σ~y~) and its value (\~ 20 Pa at 12 mg/mL) is sufficiently high to provide adequate mechanical stability to prohibit the samples being relaxed (dropped) during its transportation or shipment. Slow release of ME from the nanogels ------------------------------------ The release pattern of the volatile pheromone in terms of its relative rates of evaporation at a particular temperature was determined both for the liquid ME alone and ME immobilized in nanogel (12 mg/mL) under identical conditions. The extent of evaporation was investigated by monitoring the weight loss from both the samples kept at four different temperatures (10, 20, 30 and 50°C) in a temperature gradient chamber. Plots of %-weight losses *vs.* the number of weeks showed that the rates of evaporation were always greater from the ME alone compared to the ME immobilized in nanogel ([Figure 5](#f5){ref-type="fig"}). The calculated slopes from the linear fit of the data at various temperatures showed that the evaporation rate of ME increased significantly with increasing temperatures ([Figure 5e](#f5){ref-type="fig"}). ME in nanogel afforded only \~2--3% weight loss after 10 weeks at 30°C ([Figure 5](#f5){ref-type="fig"}) although it showed \~50% weight loss at 50°C indicating that the weight loss occurred \~13 times slower at 30°C compared to that at 50°C. On the other hand, at \~50°C (a temperature at the peak summer in Indian subcontinents), ME alone evaporated entirely within 3 weeks and in contrast it took \~30 weeks for the near complete evaporation of ME from the nanogel. The extrapolated graphs for the nanogel in their weight loss profiles further indicate that, it would take \~7.8 years at 30°C and 12.4 years at 20°C for the complete evaporation of ME from the nanogel ([SI](#s1){ref-type="supplementary-material"}, [Figure S3](#s1){ref-type="supplementary-material"}). Thus, the sustainable nature of the pheromone in the nanogel for a significant period of time makes the nanogel particularly useful for the agricultural pest management applications. Thus, the nanogel was found to be chemically, thermally and mechanically stable. The self-assembly of individual gelators is stabilized by hydrogen bonding donor (--OH) and acceptor ( = N--) sites, van der Waals interactions through the long aliphatic hydrocarbon chains and π-π stacking through the extended aromatic π-conjugated backbones of the gelator **1**[@b35]. The pheromone molecules immobilized in the nanofiber networks were found to be particularly resistant to evaporation at room temperature which confirms its increased shelf-life under these conditions. Evidence of enhanced shelf-life of ME in the nanogel ---------------------------------------------------- As the weight loss experiments indicated a significantly reduced evaporation rate of ME from the nanogels compared to the ME alone, a thin film of the prepared ME nanogel was cast (plate 'A', [Figure 6](#f6){ref-type="fig"}) on glass plates to study the extent of shelf-life of the pheromone. Control experiments were performed with the cast film of **1** solubilized in toluene (plate 'B') and ME alone (plate 'C'). All the three plates were exposed to a guava orchard for a number of hours. The results showed that the fruit flies were attracted specifically to the plate A and to the plate C throughout the period and not to the plate B which held the gelator **1** in toluene (a non-pheromone solvent, see [Movie S1](#s1){ref-type="supplementary-material"}). This indicates that the pheromone is biologically active in the nanogel and is responsible for attracting the pest although the compound **1** itself is inactive toward the fruit flies. The accumulation of the pests increased with time around both the plates A and C. The same plates were preserved at room temperature (\~30°C) and were exposed again in the same guava orchard after 21 days. Interestingly this time the fruit flies were attracted to the plate A only which contained the nanogel and not to the plate B or C for the entire time period of observation. This indicates that the nanogel laced with ME still retained the pheromone for sustained release and has also retained significant pest attractant property due to the higher shelf-life of ME in nanogel than the pure ME alone which evaporates away lot faster. This experiment also revealed that the fruit flies upon being attracted by ME released from the nanogels, gathered in the vicinity or on top of the plate. However, they did not die during this experiment (see [Movie S1](#s1){ref-type="supplementary-material"}). This indicates that the nanogel (**1** + ME) is non-toxic in nature and thus acts only as an attractant source for the fruit flies alone. Field trials with the nanogels ------------------------------ The nanogel containing ME was then utilized for the efficient pest management of *B. dorsalis* as a bait trap along with a technique to kill the pest. The total area where the experiments conducted was 67274 sq. ft. (0.625 hectare) with the distance of 15 ft. in between two adjacent trees. The samples containing ME alone and ME in nanogel were applied in the guava orchard each with six replications in such a way that the horizontal and vertical distance between the two traps was 60 ft. and 75 ft. respectively as depicted in [Figure S4 (SI)](#s1){ref-type="supplementary-material"}. The experiment was set up during the ripening stage of guava (9^th^ August, 2009) when the environmental temperature varied within \~35--40°C and this continued till the crops were harvested (7^th^ September, 2009). The nanogels (12 mg/mL) were taken in glass vials (5 mm inner-diameter, inset of [Figure 2d](#f2){ref-type="fig"}) and placed inside plastic bottles ([Figure 7a, b](#f7){ref-type="fig"}). Two holes were made in the bottles and were filled with water upto the white arrow level as shown in [Figure 7b](#f7){ref-type="fig"} (also *c.f.* Methods for the detailed Sampling Process). Control experiments were also performed with ME by directly adding it into water in the plastic bottles as convenient to the previously existing method[@b40]. Due to the presence of pheromone, the pests were attracted towards the hanging vials containing the nanogel only. During this event a significant number of fruit flies fell into water, which followed by their drowning and eventual death. Dead pests were finally collected by filtration and the total number of killed flies were counted ([Figure 7c, d](#f7){ref-type="fig"}) and summarized against events recorded each day in [Table S1](#s1){ref-type="supplementary-material"} (see [SI](#s1){ref-type="supplementary-material"}). The number of trap catches for the control containing only ME was significantly lower in comparison to the nanogel throughout the experiment. The number of catches for the control bottle containing ME alone also attracted pests initially. However, this showed a progressive decrease in comparison to the bottle containing nanogel and from the 8^th^ day onward no fly was attracted to this bottle. This indicates that the effectiveness of ME alone towards the *B. dorsalis* without the gelator **1** is at best limited to one week in the real field environment ([Figure 7e](#f7){ref-type="fig"}). On the contrary significant counts of *B. dorsalis* catches were continually recorded almost throughout the month with peak trap catches occurring during the first three weeks as compared to the control. After harvesting all the guavas, the infestation of *B. dorsalis* was found to be remarkably lower as shown in the results after three weeks of experiment since the number of catches got reduced. However, the nanogels were still active in releasing pheromone as revealed from the weight loss experiments which were conducted upto ten weeks. The procedure is therefore reproducible and particularly attractive due to the low cost, simplicity, selectivity and efficiency of the sampling technique. This also produced high yield and damage free crops. Statistical Analysis of Variance (ANOVA) were performed on the collected data containing the number of trap catches in each day throughout the experiment (30 days) for all the replications (6 for both ME alone and ME in nanogel) of the pheromone containing samples ([Table S1](#s1){ref-type="supplementary-material"}, [SI](#s1){ref-type="supplementary-material"}). The two-factor ANOVA calculations considering the factor 1 as treatment (number of trap catches) and factor 2 as the period of observations (days), were carried out on the transformed data \[(x + 0.5)^1/2^, x = number of trap catches in each day\] to homogenize the variances ([Table S2](#s1){ref-type="supplementary-material"}, [SI](#s1){ref-type="supplementary-material"})[@b41]. The interaction between the treatments (t) and the period of observations (d) has been calculated from the data presented in [Table 1](#t1){ref-type="table"} to obtain the extent of significance of the experiments. The two-factor ANOVA calculations showed that the computed F value was larger than the tabular F value both at probability (*p*) \< 0.05 and *p* \< 0.01 level of significance for the treatments (t), period of observations (d) as well as the interaction between the two (t × d) ([Table 2](#t2){ref-type="table"}). This suggests that the difference between the treatments, the period of observations pertaining to fruit flies\' visitation and the interaction between the treatments (t1 and t2) with the period of observations is highly significant. The coefficient of variation (CV) is 19.39% (\< 20%) ([Table 2](#t2){ref-type="table"}) indicating that the degree of precision with which the treatments were compared, showed a good index of reliability of the experiments[@b42]. The least significant difference (LSD) calculated (with AgRes software) for the treatments (ME alone and ME in nanogel), the period of observations (days) and the interaction between the two (t × d) considering the mean values are all presented in the [Table 1](#t1){ref-type="table"}. LSD results clearly indicate that the treatment 2 (t2, ME in nanogel) is more efficient than the treatment 1 (t1, ME alone) in trapping fruit flies (see AgRes results in [Table S5](#s1){ref-type="supplementary-material"}, [Supplementary Information](#s1){ref-type="supplementary-material"}) as the critical difference (CD) between the treatments is 0.193 (*p* \< 0.05). This further indicates that there is significant difference between the mean for the treatment 1 (2.32) and treatment 2 (7.34). The differences between the mean values of the periods of observations (days) were estimated using LSD considering the CD = 0.749 (*p* \< 0.05) which showed a descending order of distribution among days ([Table S3](#s1){ref-type="supplementary-material"}, [SI](#s1){ref-type="supplementary-material"}). The maximum number of fruit flies catches was recorded on the first three days. As the days progressed and guava fruits were harvested fruit flies infestation in field became gradually less significant and thus the lowest number of trap catches took place from the day 23 to day 30. LSD calculations were also performed from the mean values for the comparison of the interactions between the treatments and the period of observations (t × d) considering CD = 1.059 (*p* \< 0.05) and the relevant data are recorded in [Table S4 (SI)](#s1){ref-type="supplementary-material"}. Treatment 1 (ME alone) and treatment 2 (ME in nanogel) were at par on the first three days and from day 8, the treatment 2 showed more effectiveness than the treatment 1 ([Table S4](#s1){ref-type="supplementary-material"}, Group o). Therefore, ME immobilized in nanogel acted as an effective formulation compared to ME alone for the pest management of *B. dorsalis* in the field trial studies. Discussion ========== The present invention provides a simple and effective route to a slow delivery of pheromones from a nanogel. This does not require any use of environmentally harmful and toxic chemicals such as, cross-linkers, antioxidants, antimicrobials, pesticides, volatile organics *etc*. While other existing methods pose risk of human contacts where ME is added to water in the bait, the present method avoids any direct contact with the fruit crop and the agricultural workers with pheromone. This keeps the crop absolutely clean from the chemical contaminations unlike the prevalent practice of spraying toxic pesticides on the orchard. Moreover, the aggressive techniques like destruction of fallen and infested fruits, cold or vapor-heat treatments, hot-water treatment, forced hot-air treatment *etc.* are avoided in the present approach which makes this economically viable as well. The nanogel of pheromone ME described herein is insoluble in water, which makes it superior to hydrogels and microcapsules. The nanogel does not significantly swell and shrink in presence of water, confirming that the humidity plays no effect whatsoever on this material. This pheromone nanogel is particularly useful for the pest management even in rainy season as *B. dorsalis* shows maximum activity during the rainy season resulting in the most significant yield loss of crops. It may be further noted that, use of microcapsules encounters problems associated with capsule bursts or potential leaks. These critical drawbacks have been avoided in the present method. The flexibility in using the nanogel in any season at any temperature (\<60°C) are feasible due to the oxidative, photochemical and thermal stability of the nanogel. Also the transportation of the nanogel based samples is trouble free due to their significant mechanical strength. This is a very general method and the present finding could be potentially useful in preparing new nanogel type formulations with other pests such as *Helicoverpa armigera* (Hubner) (Lepidoptera, Noctuidae), *Scirphophaga incertulas* (Walker) (Lepidoptera, Pyralidae), *Leucinodes orbonalis* (Guenee) (Lepidoptera: Pyralidae) *etc.* with suitable LMMGs and appropriate pheromones that attract these species. This strategy may also be applicable even for kairomones for the attraction of predators and parasitoids (natural enemies). Such an approach could therefore be efficiently carried forward from the laboratory to the agricultural field. In summary, this is a simple and straightforward method where a nanogel is prepared as a controlled release system to deliver fruit fly pheromone that remains active over a sustained period in an environmentally safe manner. In view of the demonstrated efficiency of this pheromone nanogel in a fruit orchard, in rainy season and also in humid, sunny peak summer months, this approach should be useful in pest control in most parts of the world. Methods ======= Materials --------- Commercial ME was obtained from Sisco Research Laboratories Pvt. Ltd., India and was used without further purification. Gelator **1** has been synthesized as reported earlier[@b35]. Scanning electron microscopy (SEM) ---------------------------------- The nanogels were carefully scooped onto the brass stubs and were freeze-dried. The samples were then coated with gold vapour and analyzed on a Quanta 200 SEM operated at 15 kV. Transmission electron microscopy (TEM) -------------------------------------- A solution of **1** + ME (1 × 10^−4^ M) was loaded onto a carbon coated copper grid (200 mesh size) and the TEM images were recorded on a TECNAI F30 instrument. GC-MS studies ------------- Samples were analyzed on a gas-chromatograph (GC 17A, Shimadzu) coupled with a mass-spectrometer (GC-MS QP 5050 Shimadzu) with an electron impact ionization at 70 eV. The GC was equipped with a WCOT fused silica stationary phase CP-Sil 5CB open tubular capillary column DB-5 ms (30 m × 0.25 mm, I.D. 0.25 μM) using helium as carrier gas (flow rate 1.3 ml/min). The temperature was programmed from 60°C (hold 5 min) to 250°C at 10°C/min. The samples were introduced at the injector temperature of 250°C. Tentative identification of each compound was achieved based on the comparison of matching of the respective mass spectrum as recorded in the Wiley 139 computer databases. Identification of ME was ensured from a comparison with the mass spectra and the GC retention data with those of the known standards. Determination of gel melting temperatures (*T*~gel~) ---------------------------------------------------- The melting temperature of the nanogels from gel-to-sol was determined by an 'inverted vial' method. A sealed glass vial containing 0.5 mL of nanogel was gradually heated (2°C/min), and the temperature at which the gel started to flow under gravity was recorded as the gel melting temperature. Each experiment was performed in duplicate to determine the reproducibility. Weight loss studies ------------------- Mini glass containers of fixed (5 mm) diameter containing 0.2 mL of nanogel samples (in three replications for each sample) were placed inside a temperature gradient chamber (NK system, model: TG-180 CCFL-5LE) and were allowed to evaporate slowly. The weights were recorded on a weekly basis. Rheological studies ------------------- An Anton Paar 100 rheometer using a cone and plate geometry (CP 25-2) with a gap distance fixed at 0.05 mm was utilized. The rheometer has a built-in computer which converts the torque measurements into either G′ (the storage modulus) and G″ (the loss modulus) in oscillatory shear experiments. The nanogel was first scooped on the plate of the rheometer and then a stress amplitude sweep experiment was performed at a constant oscillation frequency of 1 Hz for the strain range of 0.01 to 500% at 20°C. Oscillatory frequency sweep experiments were performed in the linear viscoelastic region (0.01% stain amplitude) by varying the angular frequency range (0.1--100 rad/s) to ensure that the calculated parameters correspond to an intact network structure. Sampling process ---------------- We have taken a \~25 cm long plastic bottle of \~7 cm diameter with a closed screw-cap. This is kept hanging from the branch of a tree with the help of a hook in an upside down orientation. Two circular holes (diam. \~4.5 cm) have been made on the bottle in such a way that they face up and down opposite to each other. Two holes are useful for the facile passage of fruit flies that are attracted to pheromone. Water is introduced into the bottle through the lower hole and maintained at nearly the same level. The pheromone nanogel sample is kept in a hanging vial and its opening faces downwards just a few cm above the water level. With this arrangement, even when it rains, the excess of water gets automatically drained away from the lower hole allowing it to maintain the water level. As the opening of the nanogel vial is oriented downwards, rain water cannot enter into it. Also at the end of each day, the contents from the plastic bottle may be released by opening the screw-cap. This allows collection of dead flies after the day-long exposure in the orchard. The same bottle may be reused following the above procedure throughout the season. Field studies ------------- The experiment was conducted at Amrud Odyan Ropan under Rastriya Audhanik Mission at Gosaipur Mohab, Varanasi, Uttar Pradesh, India. Total area of the orchard was 0.625 hectare (67274 sq. ft) with the distance of 15 feet in between two adjacent trees. The experiment was conducted during rainy season for a period of 30 days from 9^th^ August to 7^th^ September, 2009. We took 6 replications for each sample containing ME in nanogel and ME alone. Dead pests were counted everyday for each bottle and were averaged from six bottles under each category. Thus the data containing ME samples represents a randomized block design (RBD model)[@b42]. Statistical analysis -------------------- The data were analyzed with the AgRes software (version 7.01) ANOVA (Analysis of Variance) Package for Researchers 1994, Pascal Intl software solutions. The data entry was executed using AGDATA AgRes (version 3.01) in the data entry module for the AgRes Statistical Software 1994 (Pascal Intl software solutions). Significance is assessed at 95% (p \< 0.05) and 99% (p \< 0.01) level of confidence. F values were calculated using the method as described in <http://www.biokin.com/tools/fcrit.html>. Author Contributions ==================== S.B. conceived the idea, designed experiments and analyzed results. S.B. also mentored the entire work and provided the necessary facilities to execute the work. D.B. and S.K.S. have contributed equally to this work. D.B. has performed the field trials, made the data analyses and other laboratory experiments. S.K.S. has prepared the nanogel and studied the physical properties. All contributed to the discussion of the results and writing of the manuscript. Supplementary Material {#s1} ====================== ###### Supplementary Information Supplementary Information ###### Supplementary Information Movie S1 DB thanks the present and past Directors of NBAII, Bangalore for their encouragements. We thank DST, Government of India (J C Bose Fellowship to SB), INI, IISc for the SEM and TEM studies, DBT, Government of India for funding this work. We thank Dr. K. J. David for the identification of *Bactrocera dorsalis* (Hendel). {#f1} {#f2} {#f3} {#f4} {#f5} {#f6} {#f7} ###### Mean values of the trap catches from the six replications per treatment (ME alone and ME in nanogel)[a](#t1-fn1){ref-type="fn"} Days ME alone (t1) ME in nanogel (t2) Mean Days ME alone (t1) ME in nanogel (t2) Mean -------------------------------------------------------- ----------------------------------- ------------------------------------- ---------------- ------ --------------- -------------------- -------------- 1 105 (10.18) 120.5 (10.95) 112.75 (10.57) 16 0 (0.71) 64.33 (8.01) 32.17 (4.36) 2 101.66 (10.06) 104.33 (10.21) 103 (10.14) 17 0 (0.71) 91.33 (9.47) 45.67 (5.09) 3 105 (10.18) 109.83 (10.46) 107.42 (10.32) 18 0 (0.71) 88.83 (9.4) 44.42 (5.05) 4 67.33 (8.19) 115 (10.69) 91.17 (9.44) 19 0 (0.71) 105.5 (10.19) 52.75 (5.45) 5 47.83 (6.87) 97.66 (9.87) 72.75 (8.37) 20 0 (0.71) 74.33 (8.62) 37.17 (4.66) 6 23.66 (4.81) 100 (9.99) 61.83 (7.4) 21 0 (0.71) 32 (5.42) 16 (3.06) 7 10.33 (2.96) 89.5 (9.45) 49.92 (6.20) 22 0 (0.71) 8.17 (2.69) 4.08 (1.7) 8 0.33 (0.85) 99.16 (9.94) 49.75 (5.39) 23 0 (0.71) 7.17 (2.43) 3.58 (1.57) 9 0 (0.71) 98.83 (9.89) 49.42 (5.3) 24 0 (0.71) 5.83 (2.16) 2.92 (1.43) 10 0 (0.71) 103.33 (10.15) 51.67 (5.43) 25 0 (0.71) 7.66 (2.34) 3.83 (1.52) 11 0 (0.71) 105.83 (10.24) 52.92 (5.47) 26 0 (0.71) 15.16 (3.69) 7.58 (2.20) 12 0 (0.71) 82.16 (9.07) 41.08 (4.89) 27 0 (0.71) 5.83 (2.21) 2.92 (1.46) 13 0 (0.71) 95.5 (9.75) 47.75 (5.23) 28 0 (0.71) 2.5 (1.54) 1.25 (1.12) 14 0 (0.71) 98.33 (9.91) 49.17 (5.31) 29 0 (0.71) 1 (1.09) 0.5 (0.9) 15 0 (0.71) 87.5 (9.32) 43.75 (5.02) 30 0 (0.71) 1 (1.13) 0.5 (0.92) Grand Mean 15.37 (2.32) 67.27 (7.34) 41.32 (4.83) Critical Difference (CD) between each treatments (t) 0.193[\*](#t1-fn1){ref-type="fn"} 0.255[\*\*](#t1-fn1){ref-type="fn"} Critical Difference between periods of observation (d) 0.749[\*](#t1-fn1){ref-type="fn"} 0.987[\*\*](#t1-fn1){ref-type="fn"} Critical Difference between the interactions (t × d) 1.059[\*](#t1-fn1){ref-type="fn"} 1.396[\*\*](#t1-fn1){ref-type="fn"}  ^a^This was analyzed from the data shown in [Table S1](#s1){ref-type="supplementary-material"} and [Table S2](#s1){ref-type="supplementary-material"} in the [supplementary information](#s1){ref-type="supplementary-material"}. Figures in the parentheses are the mean values for the transformed data. \[\**p* \< 0.05, \*\**p* \< 0.01\]. ###### Two-factor ANOVA calculations of the transformed data of guava field trials             Tabular F value ---------------------- ----- --------- --------- -------- ------------------------------------- ----------------- ------ Treatment (t) 1 2268.84 2268.84 2610 0.000[\*\*](#t2-fn1){ref-type="fn"} 3.87 6.72 Days (d) 29 2937.61 101.29 116.53 0.000[\*\*](#t2-fn1){ref-type="fn"} 1.51 1.78 Interactions (t × d) 29 1171.92 40.41 46.49 0.000[\*\*](#t2-fn1){ref-type="fn"} 1.51 1.78 Replications 5 1.45 0.291 0.33   2.24 3.08 Experimental error 295 256.44 0.87 1       Total 359 6636.26 18.48 21.26       \*\*Significant at *p* \< 0.05 and *p* \< 0.01; Coefficient of variation (CV) = 19.39% \[(CV) = {(Error Mean square)^1/2^/Grand Mean} × 100\]

Subscribe To Sunday, October 30, 2011 What is fructose? Fruit sugar, right? Well, yes and no. It is found in copious amounts in fruit, of course, but so are other sugars. Free fructose, the monosaccharide form, is 57% of the total sugars found in an average apple, but free glucose, another monosaccharide, is 23% of the total sugar. Sucrose, a disaccharide composed of equal parts fructose and glucose, is the remaining 20% of the sugar. So, combining the free fructose with the fructose bound up in sucrose, the total fructose in an apple is 67% of the sugars. (Trust me on the math here.) The remaining one-third is glucose. Apples and pears are on the high end of the fructose scale. Apricots, at 39%, are at the low end. The sugar in bananas is 50% fructose, grapes 53%, and peaches 46%. Honey is 50.5% fructose (free and combined). Besides tree and vine fruits, fructose is also found in other foods found in nature, for example, berries, sweet corn and sweet red peppers and most root vegetables (e.g., red beets, carrots, onions and sweet potatoes).Generally, most of the fructose is bound up in sucrose (equal parts fructose and glucose). Sucrose in its processed form is what we know as table sugar, which is made from refining sugar cane or sugar beets. Table sugar is therefore 50% fructose. So is the sucrose, or simply “sugar,” listed near the top in the ingredients list of more and more processed foods. According to Wikipedia, “Commercially, fructose is usually derived from sugar cane, sugar beets and corn, and there are 3 commercially important forms:” 1) processed crystalline fructose, 2) high-fructose corn syrup (HFCS), as a mixture of both glucose and fructose as monosaccharides, and 3) sucrose. All forms of fructose, according to Wiki, are commonly added to foods and drinks for palatability, taste enhancement and improved browning of foods such as baked goods. Starting in the early 70’s, as total consumption of sugar rose in the U.S., HFCS eroded the sucrose market. By 2000 they were consumed in the U.S. in equal amounts. HFCS is commonly found in food and drink in two forms: The 55% fructose/41% glucose form is in use in the U.S. in non-dietary soft drinks. The 42% fructose/53% glucose formulation is used primarily in processed foods and baked goods. (The balances in both forms are “other sugars.”) “The primary reason fructose is used commercially in foods and beverages, besides its low cost, is its high relative sweetness. It is the sweetest of all naturally occurring carbohydrates; at room temperature it is 1.73 times as sweet as sucrose,” but when heated it loses this advantage, again according to Wiki. The sweetness of fructose is “perceived earlier,” has a “higher peak,” “exhibits a synergy effect when used in combination with other sweeteners,” has “greater solubility,” “increases starch viscosity more rapidly and achieves a higher final viscosity than sucrose,” “retains moisture for a long period of time even at low relative humidity,” and therefore “can contribute to improved quality, better texture, and longer shelf life to the food products in which it is used,” all as reported in the Wiki entry. Ever wonder why a Twinkie or a Devil Dog stays soft forever? It’s the HFCS! If you haven’t noticed how ubiquitous HFCS has become in the processed food supply, let me give you a snapshot. In the bread aisle at my local supermarket I found it in most of the “soft” goods and long shelf life items: Devil Dogs and Twinkies, of course, and fruit pies and muffins; also in hot dog and hamburger rolls and, naturally, in Wonder Bread. I also found it listed as 4th ingredient in Weight Watchers 100% Whole Wheat bread, just before molasses! Fundamentally, however, regardless of whether the formulation of fructose you consume is 55%, 42%, or 50% fructose as in table sugar (sucrose), we all consume ever increasing amounts of fructose each year, whether we know it or not. We eat much more fructose than we think, and much more than the amount that is found in fresh fruit. Sugar -- ordinary table sugar, made from sugar cane -- remember, is half fructose. So, why does it matter? Because fructose, in the words of Robert H. Lustig, MD, is “poison.” Dr. Lustig is professor of Clinical Pediatrics in the Division of Endocrinology at the University of California San Francisco. His research focuses on childhood obesity. He contends that, in the amounts we are eating it, fructose is toxic to the liver. Want to know why? You can watch his 90 minute 2009 video, “Sugar: The Bitter Truth,” from UCSF’s “Mini Med School for the Public” on YouTube. Or stay tuned. In the next installment I will present my ‘Executive Summary’ of his evidence. Sunday, October 23, 2011 I ran into a friend at the supermarket the other day and asked her how she was. She replied, “Fine,” but “tired a lot.” Not wanting to miss an opportunity to proselytize, I suggested, “sugar crash.” She protested, “I don’t eat sugar!” So, I asked her what she ate for breakfast. Therein lies a tale of folly that deserves to be more widely known and understood. Sugar, as she and virtually the “whole world” thinks of it, is table sugar, as in “added sugar” such as that sprinkled on cereal. It is also an ingredient in candy and ice cream. Soft drinks usually are sweetened with high fructose corn syrup. Table sugar, the “added” sugar, is cane sugar and chemically is known as sucrose. It is a disaccharide, meaning it is composed of two simple sugar molecules. One of those molecules is fructose and the other is glucose. All chemical compounds ending in “ose” are sugars. Most break down in the digestion process to glucose, some to fructose, and a few to galactose, the third monosaccharide, before entering the blood stream through the wall of the small intestine. Sugars in this “whole-world” sense do not include those found in fruits, which are a combination of free fructose, free glucose and sucrose. Fruit sugars are regarded in this “whole-world” view as good for you, because they are an inherent component of this “real” or whole food. Forget that for centuries hybridizers have been making fruits sweeter than those found in nature to appeal to our sweet tooth. Carbohydrates, one of the three basic elements of nutrition (the other two being protein and fat), are all saccharides. Carbohydrates, including fruits, cereals, bread, potatoes, rice and pasta, are somewhat more complex compounds, meaning composed of many molecules. Nevertheless, they virtually all break down in the digestion process to the simple sugars glucose and fructose. Glucose goes to the cells for quick energy. Fructose goes directly to the liver and is stored. Back to the question I asked of my friend: “What did you eat for breakfast?” Her answer: “A glass of orange juice, a whole grain cereal ‘with 3 grams of protein’ [in reduced-fat milk, I assumed], toast and jelly.” “All sugar,” I exclaimed! Obnoxiously, I said, “if you eat a lot of ‘sugar’ – read any food except fat and protein -- all at once, your blood stream will be flooded with sugar (glucose) sooner or later (up to a few hours, depending on the food and the condition of your metabolism), and insulin pumped from the pancreas (if it is still working well). And then, after the ‘sugar’ gets delivered to the muscles and organs by the insulin, your blood sugar level will crash and you will ‘feel tired’ (and hungry) again.” You will feel hungry again because another hormone, ghrelin (first reported in Nature in 1999!), will send a signal from cells lining the fundus of the stomach to the hypothalamus in the center of your brain, that the “quick energy” in your blood (glucose) is low again. So, when your “all sugar” breakfast is digested, in mid-morning your body will again crave “sugar” (anything that will break down to glucose) to “feed the beast.” “Sugary snack” doesn’t mean a candy bar. It means any carbohydrate, including fruit, or a glass of milk (lactose), all of which will break down to glucose and again raise the level of sugar in the blood. It will also overwork the pancreas again to produce more insulin. A vicious cycle. Over a course of years the cell wall of the destination cells in many people will develop insulin resistance, requiring more insulin to get the job of delivering glucose energy to our muscles done. Eventually, in many of these people, the pancreas will slowly burn out. The islets of langerhans that produce the beta cells in the pancreas will stop working. They will die. By the time your doctor discovers this, up to 80% of your pancreatic function will probably already have been lost. That was the stunning conjecture made by Dr. Ralph DeFronzo, American Diabetes Association keynote speaker at their annual meeting in San Francisco in 2008. You will be diagnosed with full-blown Type 2 Diabetes. You will be drug dependent for the rest of your life. You will then be watchful for, or worse, diagnosed with the “dreaded complications.” As a result, in recent years the standards and methodology for diagnosing Metabolic Syndrome and Type 2 diabetes has evolved. The old standard was two consecutive elevated plasma blood glucose tests above 140mg/dl, lowered to 126 in 1997, where it remains today. Fasting blood glucose between 100 and 125mg/dl is now regarded as pre-diabetes. The new diagnostic standard for T2 diabetes is the Hemoglobin (Hb) A1c test. The diagnostic standard was formerly 7.0%. A few years ago it was lowered to 6.5%, and a 6.0% level added for pre-diabetes. Some activist endocrinologists, such as Dr. Richard K. Bernstein, use a much lower standard, regarding an A1c of 5.8% as indicating full-blown Type 2. Sunday, October 9, 2011 The contemporary medical literature is replete with macro and other large studies that attempt to extrapolate a correlation between heart attack risk and blood lipids. The last two columns, “Understanding Your Lipid Panel” and “The Cause and Treatment of Heart Disease” address this issue from different directions. In this column we will attempt to put a fine point on the critical matter of lipid ratios. In the 1960’s Total Cholesterol (TC) became a common and inexpensive test. As the principal metric of “the Lipid Hypothesis” popularized by Ancel Keys and then the American Heart Association, it became the universal marker for predicting heart disease risk. But that was more than half a century ago. We’ve come a long way since then. The “treatment,” i. e. medical advice, then and now for high TC was to eat less saturated fat and other animal foods with high cholesterol content. This was the modality even though the body needs cholesterol for many essential purposes and makes up what we don’t eat by manufacturing it as needed. The threshold for TC was and remains today 200mg/dl. Low density lipoprotein (LDL) is a component of TC. Even though the common test used to determine LDL was and is a calculated value, not a direct measurement, it became a popular target in the 80’s when big pharma developed drugs – statins – that lowered it. By lowering LDL, statins also lowered total cholesterol. So, doctors prescribed statins to anyone and everyone whose TC was over 200. Today, Lipitor, Crestor, Zocor, and its generic Simvastatin, account for $20 billion in world-wide annual sales. In recent decades the other components of the Lipid Panel – High Density Lipoproteins (HDL) and Triglycerides (TG) – have taken on increased importance in understanding Cardio Vascular Disease (CVD) risk.Unfortunately, these developments have garnered little attention since the pharmaceutical industry has not yet developed blockbuster drugs to influence them. Fish oil lowers Triglycerides, but fish oil cannot be patented. Most lipid panel lab results these days do however include a ratio of TC to HDL with a recommendation that it should be less than 5.0. In other words, if TC is 200, then HDL should not be less than 40. While this at least recognizes the importance of HDL, it is hardly a standard to be emulated. It is, in fact, borderline dangerous. Optimal is ≤3.5. A somewhat higher standard coming into wider use is the inverse of this fraction, i.e. HDL/TC. However, the standard for this fraction is ≥0.24 = ideal. Translated, that is closer to a ratio of 4.0, versus 5.0 in the TC to HDL ratio cited above. Many enlightened practitioners today, however, use the ratio of Triglycerides to HDL (TG/HDL) as “the single most powerful predictor of extensive coronary heart disease among all the lipid variables examined,” according to just one of many articles in the literature. The study I quote is in Clinics at PubMed Central 2008 August 63(4) 427-432. Note, by the way, that neither TC nor LDL is a factor in this formula. This ratio is considered by informed clinicians today as more reliable than LDL, or TC/HDL, or high sensitivity (hs) C-reactive protein (CRP), the marker my internist/cardiologist uses. Using this new gold standard, a TG/HDL ≤ 1.0 is considered ideal, a ratio of ≤2.0 is good, a ratio of 4.0 is considered high and 6.0 much too high. My recent TG/HDL = 0.35, which is interpreted to mean a very low probability of heart attack. As the patent on Lipitor is about to expire, and the other name-brand statin drug patent expirations are not too far behind, big pharma is hard at work looking for the next blockbuster drug to lower Triglycerides or raise HDL. Alas, so far, diet -- that is, our dietary intake (as in, “we are what we eat”) -- is the only thing that seems to work, and big pharma isn’t in that business. Agribusiness, however, has seen the potential for a big piece of the action here. Unfortunately, there isn’t much profit in “real food.” Ask your local farmer. The increasingly popular processed foods -- the so-called “heart healthy” foods we are being encouraged to eat in large quantities -- do not improve the TG/HDL ratio. They make it worse! In the coming weeks we will return to the subject of healthy eating. Subjects will include “Sugar: What do we mean?” “Fructose: Where and what is it?,” “Intermittent Fasting: Is it a good idea?,” “Ketosis and Autophagy,” and “Cooking with Oils: good and bad choices,” Our goal will be to help the reader improve their TG/HDL ratio. Of course, to do that, you will need to have a baseline Lipid Profile. If you haven’t had a Lipid Panel done, or don’t know yours, ask your doctor to do one. And ask him to see how yours shapes up using the new gold standard for CVD risk: the TG/HDL ratio. About Me I was diagnosed a Type 2 diabetic in 1986. I started a Very Low Carb diet (Atkins Induction) in 2002 to lose weight. I didn’t realize at the time that it would put my diabetes in clinical remission, or that I would be able to give up almost all of my oral diabetes meds. I also didn’t understand that, as I lost weight and continued to eat Very Low Carb, my blood lipids would dramatically improve (doubling my HDL and cutting my triglycerides by 2/3rds) and that my blood pressure would drop from 130/90 to 110/70 on the same meds. Over the years I changed from Atkins to the Bernstein Diet (designed for diabetics) and, altogether lost 170 pounds. I later regained some and then lost some. As long as I eat Very Low Carb, I am not hungry and I have lots of energy. And I no longer have any of the indications of Metabolic Syndrome. My goal, as long as I have excess body fat, is to remain continuously in a ketogenic state, both for blood glucose regulation and continued weight loss. I expect that this regimen will continue to provide the benefits of reduced systemic inflammation, improved blood lipids and lower blood pressure as well.

Factor VIII (FVIII) is a protein found in blood plasma that acts as a cofactor in the cascade of reactions leading to blood coagulation. Hemophilia A is caused by a reduction or deficiency of functional FVIII protein and is one of the most common bleeding disorders that affects about 1 in 5000-10000 men. Clinical symptoms in hemophilia are frequent muscle and joint bleeds, and trauma can even lead to life threatening situations. Currently, effective treatments for hemophilia include replacing the missing FVIII protein using intravenous application of recombinant or plasma derived FVIII products. Such preparations are generally administered either in response to a bleeding episode (on-demand therapy) or at frequent, regular intervals to prevent uncontrolled bleeding (prophylaxis). Unfortunately, the appearance of neutralizing anti-FVIII antibodies (FVIII inhibitors) is a major complication during replacement therapy with FVIII products. Approximately 25% of the patients receiving treatment develop this immunity to FVIII protein, thus making further control of bleeding very difficult. The cause for this immune response to FVIII protein has not been fully elucidated, but the specifics of a patient's immune system can affect their response to therapy. Normally, the immune system develops a tolerance to certain antigens, e.g., “self” antigens. This feature is important because, otherwise, if a self antigen is recognized as a foreign antigen, autoimmune disease results. Hemophilia A patients, in particular, have a genetic defect in their FVIII gene, which causes the immune system to not recognize the administered FVIII protein as a “self” antigen. Thus, when FVIII protein is administered during coagulation factor replacement therapy, the patient's immune system recognizes the FVIII protein as a foreign antigen or an altered self protein and develops anti-FVIII antibodies accordingly. The FVIII inhibitors, i.e., anti-FVIII antibodies are produced by plasma cells derived from FVIII specific B cells. B cells need the help of activated CD4+ T-cells to proliferate and differentiate into anti-FVIII antibody producing plasma cells. For example, FVIII protein is recognized by B and T lymphocytes in different ways. The induction of anti-FVIII antibodies is T helper cell dependent. B cells recognize whole protein epitopes via their specific B cell receptor. T-cells on the other hand, recognize proteins in the form of processed peptides complexed with an MHC class II molecule presented on the surface of an antigen presenting cell. Each CD4+ T-cell clone recognizes only one specific peptide-MHC complex. For presenting the peptides to the T-cells, MHC class II molecules have an open binding groove that allows peptides of various lengths to fit in and be presented on the surface of a cell. Moreover, the MHC class II protein contains four binding pockets that differ for the various haplotypes (Jones et al., Nature Rev. Immunol. 6:271-282 (2006)). Only specific amino acids fit into these binding pockets, and the minimal size of binding peptides is nine amino acids. Notably, different MHC class II haplotypes can present different peptides. Thus, it is likely that a patient's MHC class II haplotype influences the risk of developing anti-FVIII antibodies. Indeed, several studies have shown that there is a correlation of the human MHC class II haplotype HLA-DRB1*1501 with an increased risk for anti-FVIII antibody development (Pavlova et al., J. Thromb. Haemost. 7:2006-2015 (2009); Oldenburg et al., Thromb. Haemost. 77:238-242 (1997); Hay et al., Thromb. Haemost. 77:234-237 (1997)). Certain approaches have been explored to address the challenges associated with treating hemophilia by administration of FVIII protein. For example, WO 03/087161 discloses modified FVIII proteins, in which the immune characteristics of the FVIII protein are modified by reducing or removing the number of potential T-cell epitopes present on the protein. A number of regions that include T-cell epitopes along the FVIII protein were identified, including, e.g., FVIII2030-2044. According to the disclosure, removal of such regions could be used to provide functional FVIII protein that did not induce production of anti-FVIII antibodies. WO 09/071,886 also discloses specific regions of FVIII protein that were predicted to give rise to HLA-DR2 binding peptides that are involved in a patient's immune response, such as, e.g., FVIII475-495 FVIII542-562, FVIII1785-1805, and FVIII2158-2178. The peptides were identified for possible use in inducing immune tolerance in a patient. While there have been advances in identifying regions of FVIII protein involved in the immune response, there is still a need to identify other regions of FVIII protein that can be used for developing other therapeutic peptides and methodologies that can, for example, be used to treat patients having hemophilia A.

The practical limits of optical lithography assume that the medium through which imaging is occurring is air. This practical limit is defined by the effective wavelength equation Λ eff = λ 2 · n · NA ,where λ is the wavelength of incident light, NA is the numerical aperture of the projection optical system, and n is the index of refraction of the medium. Now, by introducing a liquid (instead of the air) between a last lens element of the projection optical system and a wafer being imaged, the refractive index changes (increases), thereby enabling enhanced resolution by lowering the effective wavelength of the light source. Lowering a light source's wavelength automatically enables finer resolution of smaller details. In this way, immersion lithography becomes attractive by, for instance, effectively lowering a 157 nm light source to a 115 nm wavelength, thereby gaining resolution while enabling the printing of critical layers with the same photolithographic tools that the industry is accustomed to using today. Similarly, immersion lithography can push 193 nm lithography down to 145 nm. In theory, older technology such as the 193 nm tools can now still be used. Also, in theory, many difficulties of 157 nm lithography—large amounts of CaF2, hard pellicles, a nitrogen purge, etc. —can be avoided. However, despite the promise of immersion photolithography, a number of problems remain, which have so far precluded commercialization of immersion photolithographic systems. These problems include optical distortions. For example, during immersion lithography scanning, sufficient g-loads are created that can interfere with system performance. These accelerative loads can cause a vibrational, fluidic shearing interaction with the lens resulting in optical degradation. The up and down scanning motions within the lens-fluid environment of Immersion Lithography can generate varying fluidic shear forces on the optics. This can cause lens vibrational instability, which may lead to optical “fading”. Other velocity profile non-uniformities can also cause optical distortions.

'''Carving''' is the practice of searching an input for files or other kinds of objects based on content, rather than on metadata. File carving is a powerful tool for recovering files and fragments of files when directory entries are corrupt or missing, as may be the case with old files that have been deleted or when performing an analysis on damaged media. Memory carving is a useful tool for analyzing physical and virtual memory dumps when the memory structures are unknown or have been overwritten. '''Carving''' is the practice of searching an input for files or other kinds of objects based on content, rather than on metadata. File carving is a powerful tool for recovering files and fragments of files when directory entries are corrupt or missing, as may be the case with old files that have been deleted or when performing an analysis on damaged media. Memory carving is a useful tool for analyzing physical and virtual memory dumps when the memory structures are unknown or have been overwritten. − based on the input's content. but it's possible (and sometimes practical) to carve individual files or [[physical memory]]. =File Carving= =File Carving= Revision as of 20:59, 29 December 2006 Carving is the practice of searching an input for files or other kinds of objects based on content, rather than on metadata. File carving is a powerful tool for recovering files and fragments of files when directory entries are corrupt or missing, as may be the case with old files that have been deleted or when performing an analysis on damaged media. Memory carving is a useful tool for analyzing physical and virtual memory dumps when the memory structures are unknown or have been overwritten. File Carving Most file carvers operate by looking for file headers and/or footers, and then "carving out" the blocks between these two boundaries. Semantic Carving performs carving based on an analysis of the contents of the proposed files. File carving should be done on a disk image, rather than on the original disk. Many carving programs have an option to only look at or near sector boundaries where headers are found. However, searching the entire input can find files that have been embedded into other files, such as JPEGs being embedded into MicrosoftWord documents. This may be considered an advantage or a disadvantage, depending on the circumstances. Today most file carving programs will only recover files that are contiguous on the media. DFRWS2006 featured a file carving challenge. As a condition of entering the challenge, all tools and techniques developed to solve the challenge had to be open sourced.

There have been proposed a variety of photovoltaic elements such as solar cells and power sources for commercial and home appliances. They utilize pn junctions formed by ion implantation or thermal diffusion of impurities into a single crystal substrate of silicon (Si) or gallium arsenide (GaAs), or by epitaxial growth of an impurity-doped layer on such single crystal substrate. However, there is a disadvantage for any of these photo voltaic elements that they are still costly because of using an expensive specific single crystal substrate. Hence, they have not yet come into general use as solar cells or power sources in commercial and home appliances used by the general public. In order to solve this problem, there have been proposed a photovoltaic element in which there is utilized a pin junction formed by amorphous silicon (hereinafter referred to as "a-Si") semiconductor films laminated on an inexpensive substrate of a non-single crystal material such as glass, metal, ceramic, synthetic resin, etc. by way of a glow discharge decomposition method. This photovoltaic element does not provide a photoelectric conversion efficiency as high as that provided by the foregoing pn junction photovoltaic element in which a single crystal substrate is used. However, this photovoltaic element can be relatively easily produced and is of low production cost, and because of this, it is used as a power source in some kinds of appliances with very small power consumption such as electronic calculators and wrist watches. In this pin junction amorphous silicon photovoltaic element, the Fermi level of the a-Si semiconductor having a good photoconductive property lies a little toward the conduction band from the center of the band gap and the electric field strength at the interface of the p-i junction is greater than that at the interface of the n-i junction. In this respect, it is advantageous to impinge light from the side of the p-type semiconductor layer in order to provide a desirable photoelectric conversion efficiency. For the p-type semiconductor layer, it is desired to be formed of such a semiconductor film that does not absorb light and does not have defects since the light to be absorbed within the p-type semiconductor layer does not contribute to generation of photoelectric current in the case where defects acting as recombination centers are present therein. In view of this, for the semiconductor film to constitute the p-type semiconductor layer in the pin junction a-Si photovoltaic element, studies have been made on amorphous silicon carbide films (hereafter referred to as "a-SiC film") which are of wide band gap and also on microcrystal line silicon films (hereinafter referred to as ".mu.C-Si film") which are known as indirect semiconductor films having small absorption coefficients and which are considered to hardly absorb light when they are of 100 to 200 .ANG. in thickness even in the case where they are of narrow band gap. As for the a-SiC semiconductor film, there is an advantage that its band gap can be widened by increasing the composition ratio of the constituent carbon atoms. However, there is a disadvantage that when its band gap is more than 2.1 eV, its film quality is markedly worsened. Therefore, there is a limit for the a-SiC semiconductor film to be used as the p-type semiconductor layer in a pin heterojunction photovoltaic element. As for the uC-Si semiconductor film, there is still a disadvantage that its band gap is narrow in any case and the quantity of light absorbed thereby is remarkable. Particularly, when the incident light is such that it contains short-wavelength light in a large proportion, the quantity of light absorbed becomes great. In view of this, in order to provide a desirable pin heterojunction photovoltaic element of the type wherein light is impinged from the side of the p-type semiconductor layer, it necessitates the use of a p-type semiconductor film having a desirably wide band gap and a minimized defect density as the p-type semiconductor layer. The same situation is present also in the case of a pin heterojunction photovoltaic element of the type wherein light is impinged from the side of the n-type semiconductor layer. That is, the n-type semiconductor layer is required to be constituted by such an n-type semiconductor film having a desirably wide band gap and a minimized defect density. Further, in the case of a so-called tandem stacked type photovoltaic element or a triple cell tandem stacked type photovoltaic element comprising a plurality of stacked cells being stacked, each cell of which comprises a pin heterojunction photovoltaic element in which the residual components of light which are left not absorbed by the upper cell are absorbed by the lower cell to obtain a sufficient photoelectric conversion, both the p-type semiconductor layer and the n-type semiconductor layer of each of the cells are required to have a desirably wide band gap and a minimized defect density. Further, for any of the foregoing photovoltaic elements, it is required for the material to constitute the p-type or n-type semiconductor layer to be such that it can be directly deposited on a non-single crystal substrate of glass, metal, ceramic or synthetic resin in a desired state and does not give any negative effect to the i-type semiconductor layer laminated thereon. As semiconductor films capable of providing a wide band which satisfy the foregoing requirements, BP semiconductor films have been proposed by Japanese Patent Laid-open No. 116673/1981 (called "literature 1" hereinafter), Japanese Patent Laid-open No. 189629/1986 (called "literature 2" hereinafter) and Japanese Patent Laid-open No. 189630/1986 (called "literature 3" hereinafter). That is, literature 1 mentions a pin heterojunction solar cell in which either the p-type or n-type semiconductor layer is comprised of a p-type or n-type amorphous BP semiconductor film (that is, a-BP semiconductor film) prepared by the glow discharge decomposition method and the i-type semiconductor layer is comprised of a a-Si semiconductor film containing fluorine atoms (F). Literature 1 does not mention anything about a crystalline BP semiconductor film (that is, poly-BP semiconductor film) which is to be distinguished from said a-BP semiconductor film. In addition, literature 1 does not describe anything about the characteristics required for a solar cell for the said pin heterojunction photovoltaic element. Literatures 2 and 3 are concerned with methods of forming semiconductor films containing group III-V elements of the Periodic Table by way of the HR-CVD method (Hydrogen Radical Assisted CVD method). But none of literatures 3 and 4 mentions anything about BP semiconductor films. Further, none of literatures 1 to 3 mentions a tandem type photovoltaic element or a triple cell tandem stacked type photovoltaic element. Against this background, there is an increased social demand to provide an inexpensive photovoltaic element which exhibits a high photoelectric conversion efficiency particularly for short wavelength light and which is practically usable as a solar cell and also as a power source in various appliances.

The family of light-harvesting-related proteins (LHCs, ELIPs, HLIPs): was the harvesting of light their primary function? Light-harvesting complex proteins (LHCs) and early light-induced proteins (ELIPs) are essential pigment-binding components of the thylakoid membrane and are encoded by one of the largest and most complex higher plant gene families. The functional diversification of these proteins corresponded to the transition from extrinsic (phycobilisome-based) to intrinsic (LHC-based) light-harvesting antenna systems during the evolution of chloroplasts from cyanobacteria, yet the functional basis of this diversification has been elusive. Here, we propose that the original function of LHCs and ELIPs was not to collect light and to transfer its energy content to the reaction centers but to disperse the absorbed energy of light in the form of heat or fluorescence. These energy-dispersing proteins are believed to have originated in cyanobacteria as one-helix, highly light-inducible proteins (HLIPs) that later acquired four helices through two successive gene duplication steps. We suggest that the ELIPs arose first in this succession, with a primary function in energy dispersion for protection of photosynthetic pigments from photo-oxidation. We consider the LHC I and II families as more recent and very successful evolutionary additions to this family that ultimately attained a new function, thereby replacing the ancestral extrinsic light-harvesting system. Our model accounts for the non-photochemical quenching role recently shown for higher plant psbS proteins.

Academic and Social Benefits of a Co-enrollment Model of Inclusive Education for Deaf and Hard-of-Hearing Children. Deaf and hard-of-hearing (d/hh) students are traditionally educated within self-contained programs at residential or special day schools, within self-contained or resource classrooms in public schools, or within regular education classrooms with support provided by an itinerant teacher. The co-enrollment model offers a promising alternative in which these students are educated within a regular education classroom composed of both d/hh and hearing students and team-taught by a teacher of the deaf and a regular education teacher. This article examines the development of one such program and the social and academic performance of the d/hh students within the program. Data on social interaction between d/hh and hearing classmates suggest that specific instructional strategies that promoted students' sign language development, identified d/hh students as "sign language specialists" and grouped d/hh and hearing students during academic activities resulted in increased interaction between these two groups of students. Stanford Achievement Test scores in the areas of reading vocabulary, reading comprehension, mathematical problem solving and procedures indicate that although d/hh students scored below the national normative hearing group, reading comprehension levels exceeded the national normative sample of d/hh students during both years two and three of the program. We discuss the challenges of implementing a co-enrollment program.

Q: Meaning of 'cast' here I'm learning something from programming and I came to this: While evaluating expressions, the result is automatically updated to larger data type of the operand. But if we store that result in any smaller data type it generates compile time error, due to which we need to type cast the result. How to explain cast in the end of the paragraph? In this context, I would go as throw something away, but yet I don't get it A: 'type cast' is a single idea. It's a coding/programming term meaning to specify a data type by forcing it, type-casting it - into the desired form. For instance "3" is a 'string', which is a text data type. The quotes in programming make that quite specific; it's not a number any more, it's a string, even though the string contains what to a human looks exactly like a number. Computers don't see things the way humans do. So - if we need to take that string & force it to become a number, we type-cast it to an integer [which is a whole number with no decimals] Our string "3" then becomes the integer 3 in pseudo-code... string myText = "3"; integer myNumber = (integer)myText; That (integer) in brackets is the 'forcing', the type-cast. Whether to hyphenate or not is probably a UK/US thing. I'm British, we hyphenate a lot more than the Americans. I guess it could come from 'to throw' it into a particular type, but I'm not certain. To cast a spell would be a similar use - whether that's 'thrown' might be debatable. Looking down the 10 meanings in the Oxford Dictionary for 'cast' I honestly wouldn't really know which it would fall under. It might even be more akin to casting an object from molten metal. A: Tetsujin explains the meaning in computer usage. I'd like to try to explain the etymology you're looking for. In manufacturing, if a liquid is poured into a mold and allowed to harden into a solid that is determined by the mold, it is said to be "cast in that mold". Type casting is a method of taking a variable and changing its form to a new type or taking an object of unknown type and treating it as a specific type. Just like the liquid in the manufacturing process, it is "cast into the mold" of the type and takes on that form.

Support Center Khan Learning Resources for CCSS Math Last Updated: Jul 15, 2016 08:12AM CDT Learning resources are adaptive items that are connected to standards and can be added to assessments for further instruction. Doing so is done through the curriculum maps feature. More information on this process is found on our support site and three articles are linked below. We have just created new curriculum maps that have Khan Academy learning resources attached to the standard. These resources will take students to the Khan pages that are associated with the standards on your assessment for more practice, helpful tutorial videos, and more. Here is an example from and 8th grade math standard: Positive and Negative Exponents. These new curriculum maps can be found by searching for "Common Core" in the curriculum maps tab. This search will populate a list of the different grade levels. Note that the middle grades are located at the bottom of the search. Once you select the map you want, you will need to re-assign the map to your classes. Do so from the actions drop down. Then, select the classes that you want the map to be assigned to. You can select multiple classes by holding down the ctrl key while selecting. Once the map has been re-assigned, it will be displayed on your curriculum map home screen. You now "own" that map and can edit and change it however you would like. Once you have assigned the map to your classes, you will be able to add the learning resources to your assessments. Adding Resources to Assessments: Once you have the curriculum map saved, you can add resources to your assessments. The first step is to add standards to your assessment. Once done, enable learning resources as an option for the assessment. Once you have enabled the learning resources and have standards aligned to the CCSS. You will see a list of the learning resources from the standards on your test. To include them on your assessment, simply select the ones you would like to include.

In lithographic printing, ink receptive regions, known as image areas, are generated on a hydrophilic surface. When the surface is moistened with water and ink is applied, the hydrophilic regions retain the water and repel the ink the ink receptive regions accept the ink and repel the water. The ink is then transferred to the surface of suitable materials upon which the image is to be reproduced. In some instances, the ink can be first transferred to an intermediate blanket that in turn is used to transfer the ink to the surface of the materials upon which the image is to be reproduced. Imageable elements useful to prepare lithographic (or offset) printing plates typically comprise one or more imageable layers applied over a hydrophilic surface of a substrate (or intermediate layers). The imageable layer(s) can comprise one or more radiation-sensitive components dispersed within a suitable binder. Following imaging, either the exposed regions or the non-exposed regions of the imageable layer(s) are removed by a suitable developer, revealing the underlying hydrophilic surface of the substrate. If the exposed regions are removed, the element is considered as positive-working. Conversely, if the non-exposed regions are removed, the element is considered as negative-working. In each instance, the regions of the imageable layer(s) that remain are ink-receptive, and the regions of the hydrophilic surface revealed by the developing process accept water or aqueous solutions (typically a fountain solution), and repel ink. Similarly, positive-working compositions can be used to form resist patterns in printed circuit board (PCB) production, thick-and-thin film circuits, resistors, capacitors, and inductors, multichip devices, integrated circuits, and active semiconductive devices. “Laser direct imaging” methods (LDI) have been known that directly form an offset printing plate or printing circuit board using digital data from a computer, and provide numerous advantages over the previous processes using masking photographic films. There has been considerable development in this field from more efficient lasers, improved imageable compositions and components thereof. Positive-working imageable compositions containing novolak or other phenolic polymeric binders and diazoquinone imaging components have been prevalent in the lithographic printing plate and photoresist industries for many years. Imageable compositions based on various phenolic resins and infrared radiation absorbing compounds are also well known. A wide range of thermally-sensitive compositions that are useful in thermal recording materials are described in patent GB 1,245,924 (Brinckman), whereby the solubility of any given area of the imageable layer in a given solvent can be increased by the heating of the layer by indirect exposure to a short duration high intensity visible light and/or infrared radiation transmitted or reflected from the background areas of a graphic original located in contact with the recording material. Thermally imageable, single- or multi-layer elements are also described in WO 97/39894 (Hoare et al.), WO 98/42507 (West et al.), WO 99/11458 (Ngueng et al.), U.S. Pat. No. 5,840,467 (Kitatani), U.S. Pat. No. 6,060,217 (Ngueng et al.), U.S. Pat. No. 6,060,218 (Van Damme et al.), U.S. Pat. No. 6,110,646 (Urano et al.), U.S. Pat. No. 6,117,623 (Kawauchi), U.S. Pat. No. 6,143,464 (Kawauchi), U.S. Pat. No. 6,294,311 (Shimazu et al.), U.S. Pat. No. 6,352,812 (Shimazu et al.), U.S. Pat. No. 6,593,055 (Shimazu et al.), U.S. Pat. No. 6,352,811 (Patel et al.), U.S. Pat. No. 6,358,669 (Savariar-Hauck et al.), and U.S. Pat. No. 6,528,228 (Savariar-Hauck et al.), and U.S. Patent Application Publications 2002/0081522 (Miyake et al.) and 2004/0067432 A1 (Kitson et al.). Positive-working thermally imageable elements containing thermally-sensitive polyvinyl acetals are described in U.S. Pat. Nos. 6,255,033, 6,541,181 (both Levanon et al.), U.S. Pat. No. 7,399,576 (Levanon et al.), and U.S. Pat. No. 7,544,462 (Levanon et al.), WO 04/081662 (Memetea et al.), and U.S. Patent Application Publication 2009/0004599 (Levanon et al.). Other positive-working imageable elements are described in copending and commonly assigned U.S. Patent Publication No. 2009/0162783 and U.S. Ser. No. 12/025,089 (filed Feb. 4, 2008 by Levanon et al.), U.S. Ser. No. 12/125,084 (filed May 22, 2008 by Levanon et al.), U.S. Ser. No. 12/195,468 (filed Aug. 21, 2008 by Levanon et al.), and Ser. No. 12/339,469 (filed Dec. 19, 2008 by Levanon et al.). Offset printing plates recently have been the subject of increasing performance demands with respect to imaging sensitivity (imaging speed) and image resolution as well as resistance to common printing room chemicals (chemical resistance). Often, the compositional features used to provide one desired property do not always improve other properties. While the imageable elements described in the patents, publications, and copending applications in the previous two paragraphs have provided useful advances in the art, additional improvements are still desired.

Wave energy buoys convert the energy of waves on a surface of a body of water, for example, from ocean surface waves, or waves formed in lakes, rivers, tidal estuaries, or other bodies of water, into electricity, which can be used to satisfy commercial and residential energy needs. Accordingly, wave energy buoys can, inter alia, reduce dependence on conventional fossil fuels. Conventional wave energy buoys use the rising and falling of the ocean waves to cause the buoy to move up and down, producing a resultant mechanical stroking within the buoy which is used to drive an electrical generator. Another conventional method of converting kinetic ocean wave energy into electricity is through the use of permanent magnets oscillating in relation to a conductive coil. As the magnets and coil oscillate relative to each other due to wave motion, electric current is created in the conductive coils. These conventional methods of converting kinetic ocean wave energy into electricity are generally inefficient since they only capture a small percentage of wave energy. This limitation is due in part to the device not being able to adjust to varying wave swells to maximize the generation of electricity. Thus, a conventional wave energy buoy cannot distinguish a wave having a greater height, wavelength, and/or speed, from a wave having a smaller height, wavelength, or speed. A conventional wave energy buoy with a fixed-strength magnet is prevented from maximizing the potential energy generated from waves with varying heights and force because where the magnet is too strong, the oscillation of the buoy is impeded reducing the amount of electric current generated in the coils. Further, where the fixed-strength magnet is too weak, the buoy will freely oscillate in relation to the coils, but the electric current generated in the coils will be less than the amount that could have been produced with a stronger magnetic field. Therefore, a conventional wave energy buoy using fixed-strength magnets of a constant strength magnetic field will not be able to optimize the potential electricity generated from each wave. Wave energy buoys can be organized into a wave farm, such as a coastal buoy farm. However, it is difficult to manage and control a plurality of buoys organized into a wave farm so that electricity can be efficiently generated on a large scale. The inefficiencies described herein with regard to individual conventional buoys are magnified when a plurality of buoys are networked together in a buoy farm.

Use of electrical vehicles is becoming increasingly popular due to the environmental benefits of removing pollution caused by fossil fuel burning vehicle engines from the environment, especially in densely populated urban environments. As with most mobile electrical devices, electrical vehicles carry electrical energy storage devices or batteries, which supply power to the vehicle propulsion and other systems. As can be appreciated, the vehicle batteries require recharging after usage. At present, electric vehicle recharging is accomplished by connecting a recharging connector into a vehicle receptacle manually. The vehicle receptacle is usually located in the same or similar location as a vehicle fuel fill spout. Such locations are protected by doors or lids, much in the same way vehicle fill spouts are protected from environmental elements, dirt and road debris. However, with the advent of automated electric vehicle recharging systems, the vehicle charge receptacles may be located elsewhere in the vehicle, including areas that are commonly inaccessible to users such as beneath the vehicle.

Electroplating and electropolishing of microelectronic workpieces, such as silicon wafers, typically involves immersing an electrically conductive surface of the wafer into a bath of liquid electrolyte. Electric current is passed through the electrolyte causing metal ions in the electrolyte to plate out onto the surface of the wafer, forming a plated layer or film. The electrical connection to the electrically conductive surface of the wafer may only be made in the so-called edge exclusion zone, as specified by industry standards. The edge exclusion zone is narrow, typically about 3 mm under current industry standards, and 2 mm or less in proposed future industry standards. In the past, contact rings having multiple spring-like contact fingers have been successfully used to provide the electrical connection to the wafer in the edge exclusion zone when processing standard wafers. Newer wafer level packaging (WLP) applications use a mask design providing the equivalent width of an edge exclusion zone of as little as 0.1 mm (100 microns). To make contact in such a narrow zone requires that the photoresist placement (i.e. centering on the wafer), the wafer size, the wafer centering in the contact ring, the manufacture of the contact fingers, etc. must hold a tolerance within 0.1 mm to make good electrical contact all the way around the circumference of the wafer. If the contact fingers only touch the exposed seed layer on portions of the circumference, then poor electrical contact will cause poor plating uniformity reducing device yield. It is of course very difficult to provide an electroplating apparatus where all of the contact fingers can consistently contact a 0.1 mm edge exclusion zone. In a similar WLP plating application, a wafer is sawed into individual dies which are tested to identify and discard defective dies. Then, only known good dies are placed into a molding compound layer on a substrate for further processing, essentially as a re-constituted wafer. In this way, only known good die are moved further in the manufacturing process. Use of re-constituted wafers also allows different types of dies to be placed next to each other with electrical connections between them made in subsequent processing steps. This approach is also beneficial for some manufacturing process such as multi-layer RDL (redistribution layer) Fan-Out. Generally, the molding compound layer on the re-constituted wafer results in a very small edge exclusion zone, of the order of about 0.1 mm. In addition, some manufacturing processes, for example multi-layer RDL, may involve wafers where the edge exclusion is different on each layer. These factors present engineering challenges in designing electroplating apparatus.

The second of two principles contained in the "Just War Theory," Jus Ad Bellum and Jus In Bello. Jus In Bellum in Latin means "The Law in Waging War." It defines standards by which a country can conduct war maintain that they have "just" actions in war. The term was coined in the same era in which the League of Nations originated. However, it was not used in doctrine until the late 1940s. Its principles include discrimination and proportionality. Descrimination defines legitimate targets and proportionality defines how much force to be used. Under the guidelines of Jus In Bello, the attack on Heroshimo, Japan by US forces with the use of the atomic bomb violates both the guidelines of descrimination and proportionality. A group of principles intended as guidelines for the just prosecution of war whose origins go back at least as far as the writings of St. Thomas Aquinas (ca. 1225-1274 AD). Includes the two principles of “discrimination” and “proportionality,” which attempt to define valid targets and acceptable applications of force, respectively. Real understanding of the circumstances of the attacks on Hiroshima and Nagasaki would lead one to conclude that the use of the atomic bomb was actually fully justified under the guidelines of jus in bello.

Parent-child relationships are very common in relational database schemas. For example, a database schema might indicate a relationship between a manager (the “parent”) and one or more employees (the “children”) who are subordinate to that manager. For another example, a database schema might indicate a relationship between a purchase order (the “parent”) and one or more lines (the “children”) that are a part of that purchase order. Data that contains such parent-child relationships is called hierarchical data, since the parent-child relationships form a hierarchy. Typically, the way that a parent-child relationship is modeled within a relational database table is to include, within the table that contains the “child” records, a column that indicates a row identifier for a “parent” record that is related to those child records in a parent-child relationship (in a relational database table, each row of the table may have a unique row identifier). For example, in an “employee” table, an employee “Alice” might be a child record of another employee “Mel;” Mel might be Alice's manager, for example. The employee table might contain a “manager ID” (i.e., the parent ID) column. In the employee table, the row containing Alice's record might contain, in the manager ID column, the row identifier for the row (in the employee table) that contains Mel's record. Taking the example further, if Mel is also the manager for employees “Vera” and “Flo,” then, in the employee table, the rows containing Vera's and Flo's records might also contain, in the manager ID column, the row identifier for the row that contains Mel's record. If Mel himself has a manager, then the value of the manager ID in Mel's record might also indicate the row identifier for the row that contains Mel's manager's record. Structured Query Language (SQL) has operators and constructs that allow an interested user to formulate queries that, when executed, will (a) cause all of the children of a specified parent to be returned, (b) cause all of the descendants (i.e., the children of, the children of those children, etc.) of a specified parent to be returned, (c) cause the parent of a specified child to be returned, and/or (d) cause all of the ancestors (i.e., the parent of, the parent of that parent, etc.) of a specified child to be returned. Using such an approach, determining the immediate parent or children of a specified entity may be accomplished fairly trivially using a simple index lookup. However, determining more extended relationships, such as all ancestors or all descendants of a specified entity, can be much more complicated. Although SQL contains a “CONNECT-BY” construct that permits such complicated operations to be expressed in a relatively compact manner, evaluating a query that contains such a construct involves an N-step algorithm in which each step involves a separate index scan. For example, if a query requests all of the descendants of a specified entity, and if there are N levels of descendants represented in the hierarchy, then evaluating the query will involve at least N separate index scans. Where N is very large, evaluating such a query can be very costly in terms of both time and computing resources. Database system performance can suffer when such a query is evaluated. What is needed is a way of obtaining, from a database, hierarchical relationship information, such as all of the ancestors or descendants of a specified entity, without incurring the high costs that are associated with large numbers of index scans. The approaches described in this section are approaches that could be pursued, but not necessarily approaches that have been previously conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section.

Blue Ridge Railway (1901) The Blue Ridge Railway was a 19th-century railroad in the U.S. state of South Carolina. It was originally chartered in 1852 as the Blue Ridge Railroad of South Carolina. Original plans were for a 195-mile line from Anderson, South Carolina, to Knoxville, Tennessee going through the mountains with as many as 13 tunnels including the incomplete Stumphouse Mountain Tunnel. By 1859, the railroad had built between Anderson and West Union, South Carolina and substantial work on several tunnels had started. However, with the outbreak of the Civil War, no further work was ever completed on the tunnels despite efforts after the war, including one by the Black Diamond Railroad. A final extension from West Union to Walhalla, South Carolina at the urging of the Town Council and local citizens saw the first train arrive November 14, 1877, but no more track would ever be laid along the alignment up Stumphouse Mountain. In 1880, the Columbia and Greenville Railroad acquired the of track that had been laid. In 1901 the Southern Railway, successor to the Columbia and Greenville, split out the line as the Blue Ridge Railway. The Southern leased the Blue Ridge to subsidiary Carolina and Northwestern Railway on July 1, 1951, and eventually merged the company. References External Links Blue Ridge Railroad history Blue Ridge History trail at Stumphouse Tunnel Park and Issaqueena Falls in Oconee County, South Carolina Category:Defunct South Carolina railroads Category:Predecessors of the Southern Railway (U.S.) Category:Railway companies established in 1901 Category:American companies established in 1901

Environmental Policy Lecture 8 part 1 In lecture 8 Professor Sterner speaks about discounting, an exiting topic as he calls the discount rate "the value of the future". In spite of all uncertainties there are about impacts of Climate Change the biggest uncertainty considering its future costs seems to come from the discount rate. In part 1 Sterner discusses the different concepts. He compares and explains the differences between discount rate and growth, and explains why it is impossible to have a steady high growth rate. He gives an example of the first contact he had with discount rates in the 70s in a discussion about the costs of dealing with nuclear waste. The opponents had by estimating a high discount rate argued that the future cost would be so low that it would not be a problem for future generation. This is the same logic, as Sterner shows, that Professor William Nordhaus uses in his calculations for the future costs of Climate Change, and also the main point which differs Nordhaus calculations from Professor Nicholas Stern. When estimating future values one has to be careful when assuming discount rate, Sterner calculates and shows how extremely different the outcomes are with just small differences in discount rate. "What does it mean to become 100 times richer?" A hundred years ago 10% of New York’s population had maids, now basically no one has a maid. The relative prices have gone up. Sterner further shows the importance of considering relative prices by illustrating with the market of agriculture which today is about 24% of world GDP . If agriculture would collapse and we would loose 95% it would not be the same as loosing about 23% of GDP since the still existing 5% of agriculture would basically become the entire GDP.

行, by itself, can be used as an adverb meaning "soon" or "before long," and 將 can be used as an auxiliary verb meaning "will," so 行將 translates as "will soon." The Korean equivalent is 곧. Notice that the adverb (行) comes before both the auxilliary verb (將) and the main verb (歸). Saturday, June 11, 2016 食 (식) means "to eat" and 此 (차) means "this," but what does 革 (혁) mean in the above sentence? Besides meaning "leather" or "animal hide," 革 can also mean "to reform," "to remove," or "to expel (from office)," but can it also mean "instead"? On Page 24 of "Du's Handbook of Classical Chinese Grammar," in the section explaining adverbs and modal verbs, 革食此 was one of the example sentences used to show how adverbs were placed before the verbs they modified in Classical Chinese. 革食此 (혁식차) was translated as "eat this instead," meaning that 革 was translated as "instead." Before now, I cannot remember ever seeing 革 translated as "instead," but words like 革新 (혁신), which means "reform" or "renovation," does seem to imply "instead," in that when you reform something, instead of the old way, you try a new way. Has anyone seen any other examples of 革 being used to mean "instead"?

Q: Python3--How to implement binary search in number game Hello I am working on learning binary search trees. The program explains itself. I am confused on whether or not I should write a function to half the numbers. I am just confused on the logic of it all. If you could help I would greatly appreciate it. list_of_numbers = [i for i in range(1,101)] print("THIS IS YOUR LIST OF NUMBERS", list_of_numbers) lowest_number = 1 highest_number = len(list_of_numbers)-1 mid = (len(list_of_numbers)) / 2 print("IS THIS YOUR NUMBER???:", mid) number_of_guesses = 0 while lowest_number <= highest_number: print('Is your number: ', mid) print('please enter "<,>, or ="') first_guess = input("???:::::::::") first_answer = first_guess print(first_answer) equals = "=" greaterthan = ">" lessthan = "<" if first_answer == equals: print('Great JOB I GOT YOUR NUMBER ON THE FIRST TRY') break if first_answer == lessthan: print('The number is now between 49-1, half it.') break elif first_answer == greaterthan: print('The number is now between 51-100') break else: print('Type one of the answers') break number_of_guesses = number_of_guesses + 1 A: You don't need a function to simply divide by 2. What you should be doing for a binary search is keeping track of the upper and lower bounds for the number; then you can calculate the guess as (upper + lower) / 2. If the player says the guess is too high, then set the upper bound to the current guess. If they say the guess is too low, then set the lower bound to the current guess. Repeat that until the correct number is found. BTW, this isn't really a search tree, it's just a basic binary search. A search tree would be storing all the possible numbers in a connected tree graph and using that to search for the answer, which this isn't doing.

Gas valves are commonly used in conjunction with gas-fired appliances for regulating gas flow and/or gas pressure at limits established by the manufacturer or by industry standard. In many cases, such devices can include pressure regulators to, for example, establish and/or maintain a gas pressure to help prevent over-combustion or fuel-rich combustion within the appliance, and/or to prevent combustion when the supply of gas is insufficient to permit proper operation of the appliance. Examples of gas-fired appliances that may employ such gas valves can include, but are not limited to, water heaters, furnaces, fireplace inserts, gas stoves, gas clothes dryers, gas grills, or any other such device where gas control is desired. Typically, such appliances utilize fuels such as natural gas or liquid propane as the primary fuel source, although other liquid and/or gas fuel sources may be provided depending on the type of appliance to be controlled. In a gas-fired appliance, a combustion chamber and air plenum are typically provided along with a gas valve. A burner element, fuel manifold tube, ignition source, thermocouple, and/or pilot tube can also be provided as part of the burner system. During operation, when a heat demand is present, metered fuel is typically introduced via the gas valve through the fuel manifold tube and burner element and into the combustion chamber. The fuel is ignited by a pilot flame or other ignition source, causing fuel combustion at the burner element. In some cases, air may be drawn into the air plenum, sometimes under the assistance of an air blower, causing the air to mix with the fuel to support the combustion within the combustion chamber. The products of the combusted air-fuel mixture are typically fed through a flue or heat exchanger tube in the gas-fired appliance to heat by convection and conduction. In some cases, the gas valve may include a pressure regulator to regulate the flow of gas at a pressure. In many cases, the pressure regulator references the pressure of the combustion chamber to help maintain and/or achieve a desired combustion level in the combustion chamber. Typically, a hose can be coupled to the pressure regulator and the combustion chamber or burner box to fluidly connect the pressure regulator and the reference pressure within the combustion chamber. It has been found, however, that in some cases, the hose may become blocked by condensate build-up or other particulate matter. Also, in some cases, the hose may become kinked or otherwise obstructed. In either case, the blockage may cause the pressure in the pressure regulator to increase or decrease resulting in over-combustion or under-combustion in the combustion chamber. In some gas-fired systems, a separate fitting including a bleed orifice may be coupled between the pressure regulator and the hose to provide a reference pressure such as atmosphere if the hose becomes blocked. This additional fitting can, however, be removed from the system or not installed properly during installation. Also, the bleed orifice in the fitting may become blocked with grease or other material during handling.

Roots and Radicals Roots and Radicals introduces students to rational exponents and helps with studying how to simplify and otherwise manipulate expressions with radicals and roots using their properties and rules. Students will learn through practice problems like those found in their homework how to rationalize, combine, expand radicals. Topics include:

Introduction {#s1} ============ The ability of pathogenic microorganisms to assimilate nutrients from their host environment is one of the most fundamental aspects of infection. To counteract this, hosts attempt to withhold essential micro-nutrients from potentially harmful microbes to limit, or even prevent, their growth. This process is called nutritional immunity. For example, vertebrates, such as humans, express several iron-binding molecules to maintain extremely low free levels of this metal in the body. To overcome this restriction, successful pathogens have evolved sophisticated mechanisms to assimilate iron. These include high affinity transporters, siderophores, and transferrin-, ferritin-, and haem-binding proteins [@ppat.1003034-Hood1], [@ppat.1003034-Almeida1]. Indeed, iron acquisition is considered a vital virulence factor for many pathogens. However, nutritional immunity does not begin and end with iron. Vertebrates have also developed mechanisms to sequester other essential metals, such as zinc [@ppat.1003034-Corbin1], [@ppat.1003034-Urban1]. The importance of zinc sequestration and the strategies that successful pathogens employ to overcome this has only recently been realized. Why Is Zinc So Important? {#s2} ========================= Zinc is essential for life, with an astonishing 9% of eukaryotic proteins predicted to be zinc metalloproteins [@ppat.1003034-Andreini1]. The role of this metal in human health is well-documented [@ppat.1003034-Chasapis1] and zinc is known to play key roles in both adaptive and innate immunity. However, host zinc sequestration from pathogens, as a means to control microbial growth, is an emerging field [@ppat.1003034-KehlFie1]. Accompanying this, a growing body of literature is illuminating the role of bacterial zinc acquisition systems in virulence [@ppat.1003034-Hood1]. But what about fungal pathogens? Fungi also rely on zinc for growth, as this metal serves as a cofactor for several enzymes, including superoxide dismutase and alcohol dehydrogenase, along with numerous other proteins, such as transcription factors [@ppat.1003034-Colvin1]. Therefore, in order to cause infections, pathogenic fungi must assimilate zinc from their host environment. Here, we will discuss the mechanisms of zinc exploitation by human pathogenic fungal species. How Can the Host Sequester Zinc from Potential Invaders? {#s3} ======================================================== Although zinc is the second most abundant transition metal in the human body, its spatial distribution is highly dynamic. In mammals, this is largely mediated at the cellular level by zinc transporters. These include ZIPs (Zrt-, Irt-like proteins/solute carrier family 39, SLC39), which deliver zinc into the cytoplasm and ZnTs (zinc transporter), which pump zinc out of the cell or into vesicles [@ppat.1003034-Liuzzi1]. Within the cell, zinc availability is tightly regulated via sequestration within organelles and binding to proteins such as metallothioneins [@ppat.1003034-Colvin1]. During acute inflammation, hepatocytes remove zinc from the plasma via ZIP14 [@ppat.1003034-Liuzzi2], reducing availability to extracellular pathogens. On the other hand, ZnT transporters reduce cytoplasmic zinc levels, possibly limiting access to intracellular pathogens; indeed, *Listeria monocytogenes* relies on two zinc uptake systems for intracellular growth [@ppat.1003034-Corbett1]. Furthermore, it would appear that ZIP-mediated export of zinc from phagosomes may also be employed to limit the growth of phagocytosed intracellular pathogens [@ppat.1003034-Aydemir1], [@ppat.1003034-Serafini1]. In contrast, it has also been shown that phagocytes attempt to kill intracellular pathogens, such as *Mycobacterium tuberculosis*, by increasing the heavy metal content of intracellular compartments to potentially toxic levels [@ppat.1003034-Botella1], [@ppat.1003034-Botella2]. Therefore, both metal restriction and metal overload represent potential mechanisms to control pathogenic microorganisms. In addition to cellular import/export, zinc can also be limited via calprotectin, an antimicrobial peptide with zinc (and manganese) chelation properties [@ppat.1003034-Corbin1]. Neutrophils (one of the most important antifungal effectors) contain high levels of calprotectin and these phagocytes decorate their NETs (neutrophil extracellular traps) with this potent antimicrobial peptide. Indeed, calprotectin is particularly important for the candidacidal activity of NETs [@ppat.1003034-Urban1]. Intriguingly, *Salmonella* Typhimurium can actually exploit calprotectin-mediated zinc sequestration in the inflamed guts of infected mice, thus out-competing rival commensal species [@ppat.1003034-Liu1]. Nutrient acquisition from this antimicrobial peptide represents a striking example of a pathogen gaining an upper hand in the continuous arms-race with its host. It remains an open question whether fungal pathogens can also exploit this potential zinc source during infection. How Do Fungi Obtain Zinc? {#s4} ========================= A number of studies have investigated the mechanisms of zinc homeostasis in the model yeast *Saccharomyces cerevisiae*. This fungus encodes two plasma membrane transporters, Zrt1 and Zrt2 ([Figure 1A and 1B](#ppat-1003034-g001){ref-type="fig"}), which are up-regulated by the transcription factor Zap1 and transport zinc into the cell. If cellular levels become too high, these importers are rapidly down-regulated [@ppat.1003034-Eide1]. {#ppat-1003034-g001} The fungal vacuole is also extremely important for cellular zinc homeostasis: Cot1/Zrc1 (ZnT type transporters) import zinc into the fungal vacuole, thus protecting the cell from potentially toxic levels of the metal; conversely, if the cell experiences zinc starvation, the zinc pool stored in the vacuole can be mobilized via the vacuolar zinc exporter, Zrt3 ([Figure 1B](#ppat-1003034-g001){ref-type="fig"}) [@ppat.1003034-MacDiarmid1]. Indeed, vacuolar storage and mobilization can significantly buffer the cytosol from both zinc excess and depletion. Unlike bacteria, which express cell membrane-localized zinc efflux systems [@ppat.1003034-Hantke1], fungi do not appear to actively export zinc in this manner [@ppat.1003034-Eide1] and rather rely on vacuolar sequestration to detoxify this metal. The CtpC zinc efflux system of *M. tuberculosis* is required for survival within macrophages, as these immune cells attempt to poison phagocytosed bacteria with potentially toxic levels of zinc [@ppat.1003034-Botella1], [@ppat.1003034-Botella2]. Although it is at present unknown whether phagocytosed fungi also experience zinc toxicity, vacuolar detoxification would represent a promising mechanism for counteracting this [@ppat.1003034-Yasmin1]. Indeed, orthologues of vacuolar zinc transporters are found in several pathogenic fungi. As well as counteracting metal toxicity, intra-vacuolar storage can also endow yeasts with an extraordinary capacity to withstand zinc starvation. Simm et al. have demonstrated that the storage capacity of the *S. cerevisiae* vacuole is sufficient for a single mother cell to produce many progeny, even in the absence of zinc uptake [@ppat.1003034-Simm1]. Thus vacuolar zinc storage may have serious implications for virulence studies aimed at elucidating zinc uptake systems in fungal pathogens. Zincosomes represent another potential intracellular zinc store ([Figure 1B](#ppat-1003034-g001){ref-type="fig"}). Zincosomes are vesicles that contain labile zinc and have been observed in both mammalian and yeast cells. Zincosomes may serve to both detoxify excess zinc and mobilize this metal upon deprivation (analogous to the fungal vacuole); however, the exact nature of these compartments, and the mechanisms by which they function, remain unclear [@ppat.1003034-Eide1]. Zinc acquisition by a pathogenic fungus has been most extensively investigated in *Aspergillus fumigatus* by the Calera group. *A. fumigatus* encodes three zinc transporters (*zrfA*-*C*), all of which are positively regulated by ZafA, the functional orthologue of yeast Zap1 [@ppat.1003034-Amich1], [@ppat.1003034-Moreno1]. Although the role of these transporters in *A. fumigatus* pathogenicity has not yet been directly examined, deletion of their transcriptional activator, ZafA, abrogates virulence, suggesting a role for zinc uptake during infection [@ppat.1003034-Moreno1]. Similarly, *Candida albicans* Zap1 is required for the expression of zinc transporter encoding genes and plays a crucial role in biofilm formation, an important pathogenicity attribute of this fungus [@ppat.1003034-Nobile1]. Therefore, the transcriptional regulation of zinc transporters would appear to be conserved, at least amongst *S. cerevisiae*, *C. albicans*, and *A. fumigatus*. Phylogenetic studies have identified zinc transporters in numerous other pathogenic fungi [@ppat.1003034-Citiulo1] and [Figure 1A](#ppat-1003034-g001){ref-type="fig"}; however, their roles in virulence remain largely unexplored. In addition to transporter-mediated assimilation, zinc can also be scavenged via the zincophore system. The Fungal Zincophore System {#s5} ============================ We recently reported sequestration of host zinc by the major human fungal pathogen *C. albicans* [@ppat.1003034-Citiulo1]. We hypothesized that, analogous to siderophore-mediated iron acquisition, the fungus might secrete a zinc-binding molecule to scavenge this metal from its environment---we proposed the term "zincophore" for such a scavenger. Indeed, we found that *C. albicans* secretes a zinc-binding protein (Pra1, pH-regulated antigen), which can sequester this metal from the environment and is required for stealing zinc from host cells ([Figure 1B](#ppat-1003034-g001){ref-type="fig"}). Moreover, Pra1 reassociation with the fungal cell (a prerequisite for a functional zincophore system) was found to be mediated by the plasma membrane zinc transporter (Zrt1), encoded at the same locus as the zincophore. Therefore, Zrt1 of *C. albicans* likely has dual transporter and receptor functions ([Figure 1B](#ppat-1003034-g001){ref-type="fig"}). The Calera lab has shown that the Pra1 and Zrt1 orthologues of *A. fumigatus* (Aspf2 and ZrfC, respectively) are also syntenically encoded, regulated by environmental zinc status and required for growth under zinc starvation [@ppat.1003034-Amich1], indicating functional conservation of zinc-acquisition loci between *Candida* and Aspergilli. Indeed, we found the zincophore locus in both ascomycetes and basidiomycetes, groups of fungi that diverged around half a billion years ago---clearly, this system is not specific to human pathogens. Rather, it appears to be associated with zinc foraging in certain niches of neutral/alkaline pH. In line with this, both characterized systems (in *A. fumigatus* and *C. albicans*) are strongly repressed by acidic pH, even under conditions of zinc limitation [@ppat.1003034-Amich1], [@ppat.1003034-Citiulo1]. Although the zincophore system is essential for *C. albicans* zinc scavenging during host cell invasion [@ppat.1003034-Citiulo1], many other human fungal pathogens do not encode Pra1 [@ppat.1003034-Citiulo1]. These other species must rely on alternative acquisition systems. It is possible that some species rely solely on transporters for zinc uptake. Alternatively, another, convergently evolved, secreted zinc-binding protein may be deployed. Finally, it is possible that some fungi may secrete small molecule zinc chelators to sequester this metal. The most well characterized zincophore-encoding species, *C. albicans* and *A. fumigatus*, are both notable for aggressive, hypha formation-mediated tissue invasion and inflammation. Indeed, Pra1 is immuno-modulatory [@ppat.1003034-Zipfel1] and is the major ligand for leukocyte integrin αMβ2 [@ppat.1003034-Soloviev1]. Similarly, the *A. fumigatus* homologue Aspf2 is a major allergen, which cross-reacts with over 80% of sera from patients suffering from aspergilloma or allergic bronchopulmonary aspergillosis [@ppat.1003034-Segurado1]. Therefore, an intriguing possibility is that the loss of the zincophore system by contemporary human fungal pathogens [@ppat.1003034-Citiulo1] may actually contribute to their capacity to evade certain aspects of immune recognition. Whilst foregoing zincophore-mediated zinc scavenging, yeasts such as *C. glabrata*, *Cryptococcus neoformans* and *Histoplasma capsulatum* may benefit from avoiding unwanted attention from aggressive host immune responses. In summary, as the scope of nutritional immunity expands beyond iron to encompass other metals, the molecular mechanisms that pathogenic microorganisms deploy to circumvent host metal restriction represents fertile ground for the identification of novel virulence factors. We would like to thank all of the members of MPM for ever helpful discussions and Matthias Brock and François Mayer for critical reading of the manuscript. [^1]: The authors have declared that no competing interests exist.

NASA has announced its plans for a nuclear-powered aerial drone, Dragonfly, that will fly from place to place on Saturn’s moon Titan. Dragonfly is planned to launch in 2026 for an expected landing in 2034. The lengthy flight will include a gravity-assist maneuver to increase the probe’s velocity. The wild card that may change that is the SpaceX Starship, capable of sending Dragonfly on a quicker, more direct route, which may be available as early as 2021. Dragonfly will be the first space probe that will travel on an alien world entirely by air. It will have a suite of cameras that will take pictures both at a distance and close up. Dragonfly will also carry a mass spectrometer to analyze the chemical makeup of the materials it encounters. It will be able to conduct meteorological studies while in flight and seismographic studies of Titan’s interior at each site where it lands. ADVERTISEMENT Dragonfly’s primary mission will be to search for signs of life, flying to diverse sites, including organic dunes and an impact crater where scientists believe water and prebiotic materials may have existed together in the distant past. The examination of these sites in the unique environment of Titan will advance astrobiology by tremendous leaps and bounds. Dragonfly will hop from place to place using eight helicopter rotors for distances of up to five miles. It will visit dozens of sites during the primary science mission that is scheduled to last 2.7 years. Because of Titan’s thick atmosphere and its distance from the sun, Dragonfly will be powered by a nuclear energy device similar to the one used by the Mars Curiosity rover. Titan is one of the weirdest and therefore most interesting worlds that humanity is aware of. It is the second-largest moon in the solar system. Its surface is comprised of rock and water ice. The moon’s atmosphere consists of nitrogen and methane. Liquid methane and ethane comprise lakes and streams on Titan’s surface, evaporating into clouds before condensing as rain or light snow, in short creating a weather cycle much like water does on Earth. Titan is unlikely to contain life, at least as we know it. The moon of Saturn is 886 million kilometers from the sun and regularly experiences temperatures of -290 degrees Fahrenheit, Because Titan’s atmosphere is so thick, four times that of Earth’s, the surface pressure is 50 percent higher than that of our home planet. Still, no one will know for sure about life on Titan until explorers — robotic and perhaps, in the distant future, human — traverse Saturn’s largest moon. Indeed, NASA engineer Janelle Wellons once speculated that a human settlement could be established on Titan. Titan has been extensively explored from above by the Cassini probe which orbited Saturn between 2004 and 2017. While images of Titan that were taken from space resemble an orangish blob, Cassini also brought with it the Huygens probe. That was courtesy of the European Space Agency, which entered the moon’s atmosphere and touched down on its surface after being slowed by a parachute. Huygens took images of Titan’s surface, an orange plain peppered with water ice pebbles under a misty, methane sky. The Huygens remains the most distant landing of any human spacecraft in history. Dragonfly is part of the New Frontiers series of space probes, several of which are still in operation. The OSIRIS-REx is currently in orbit around the asteroid Bennu. Juno is in a polar orbit of Jupiter and has imaged parts of the largest planet in the solar system hitherto unseen by humans. The most famous of the series, New Horizons, flew past Pluto, returning the first images of the dwarf planet and its moons. It then head out toward an unusual Kuiper Belt world called Ultima Thule. Dragonfly, it is hoped, will add to the wealth of scientific discoveries yielded by her predecessors. Mark R. Whittington is the author of space exploration studies “Why is It So Hard to Go Back to the Moon? as well as “The Moon, Mars and Beyond.”

Extra A Brief History of Cell Death Research The stereotypical developmental 'regression' of various fetal and larval structures was known to the ancients, including possibly Aristotle. By the 1840's, technical developments in microscopy and the elaboration of the cell theory by Schleiden & Schwann, allowed for such events to be placed into a cellular context by Vogt and others: developmental regression in toads and insects was understood to involve the death of individual cells. In an early example of the application of chemistry to understand biology (what today we would call 'chemical biology'), textile dyes used in the 1880's by Walther Fleming and others to stain cells subsequently enabled the visualization of new intracellular features of dying cells, such as 'chromatolysis' - the regulated destruction of the chromosomes. In the 20th Century it became clear from studies of animal development that some cells formed during early development were subsequently eliminated. The name apoptosis was coined by Kerr, Wyllie & Currie in 1972 to describe one characteristic process of programmed cell death. The nematode worm C. elegans, introduced as a genetic model of animal development in the 1970's by Sydney Brenner, was used by Robert Horvitz's lab in the 1980's and 1990's to isolate specific genes involved in this process. Converging lines of research demonstrated that cell death pathways were highly conserved and important for development and the prevention of diseases such as cancer in mammals. A subsequent explosion of work has identified numerous biochemical pathways that, when activated, can trigger some form of regulated cell death.

Évariste Galois, (born October 25, 1811, Bourg-la-Reine, near Paris, France—died May 31, 1832, Paris), French mathematician famous for his contributions to the part of higher algebra now known as group theory. His theory provided a solution to the long-standing question of determining when an algebraic equation can be solved by radicals (a solution containing square roots, cube roots, and so on but no trigonometry functions or other nonalgebraic functions). Galois was the son of Nicolas-Gabriel Galois, an important citizen in the Paris suburb of Bourg-la-Reine. In 1815, during the Hundred Days regime that followed Napoleon’s escape from Elba, his father was elected mayor. Galois was educated at home until 1823, when he entered the Collège Royal de Louis-le-Grand. There his education languished at the hands of mediocre and uninspiring teachers. But his mathematical ability blossomed when he began to study the works of his countrymen Adrien-Marie Legendre on geometry and Joseph-Louis Lagrange on algebra. Under the guidance of Louis Richard, one of his teachers at Louis-le-Grand, Galois’s further study of algebra led him to take up the question of the solution of algebraic equations. Mathematicians for a long time had used explicit formulas, involving only rational operations and extractions of roots, for the solution of equations up to degree four, but they had been defeated by equations of degree five and higher. In 1770 Lagrange took the novel but decisive step of treating the roots of an equation as objects in their own right and studying permutations (a change in an ordered arrangement) of them. In 1799 the Italian mathematician Paolo Ruffini attempted to prove the impossibility of solving the general quintic equation by radicals. Ruffini’s effort was not wholly successful, but in 1824 the Norwegian mathematician Niels Abel gave a correct proof. Galois, stimulated by Lagrange’s ideas and initially unaware of Abel’s work, began searching for the necessary and sufficient conditions under which an algebraic equation of any degree can be solved by radicals. His method was to analyze the “admissible” permutations of the roots of the equation. His key discovery, brilliant and highly imaginative, was that solvability by radicals is possible if and only if the group of automorphisms (functions that take elements of a set to other elements of the set while preserving algebraic operations) is solvable, which means essentially that the group can be broken down into simple “prime-order” constituents that always have an easily understood structure. The term solvable is used because of this connection with solvability by radicals. Thus, Galois perceived that solving equations of the quintic and beyond required a wholly different kind of treatment than that required for quadratic, cubic, and quartic equations. Although Galois used the concept of group and other associated concepts, such as coset and subgroup, he did not actually define these concepts, and he did not construct a rigorous formal theory. While still at Louis-le-Grand, Galois published one minor paper, but his life was soon overtaken by disappointment and tragedy. A memoir on the solvability of algebraic equations that he had submitted in 1829 to the French Academy of Sciences was lost by Augustin-Louis Cauchy. He failed in two attempts (1827 and 1829) to gain admission to the École Polytechnique, the leading school of French mathematics, his second attempt marred by a disastrous encounter with an oral examiner. Also in 1829 his father, after bitter clashes with conservative elements in his hometown, committed suicide. The same year, Galois enrolled as a student teacher in the less prestigious École Normale Supérieure and turned to political activism. Meanwhile he continued his research, and in the spring of 1830 he had three short articles published. At the same time, he rewrote the paper that had been lost and presented it again to the Academy—but for a second time the manuscript went astray. Jean-Baptiste-Joseph Fourier took it home but died a few weeks later, and the manuscript was never found. The July Revolution of 1830 sent the last Bourbon monarch, Charles X, into exile. But republicans were deeply disappointed when yet another king, Louis-Philippe, ascended the throne—even though he was the “Citizen King” and wore the tricoloured flag of the French Revolution. When Galois wrote a vigorous article expressing pro-republican views, he was promptly expelled from the École Normale Supérieure. Subsequently, he was arrested twice for republican activities; he was acquitted the first time but spent six months in prison on the second charge. In 1831 he presented his memoir on the theory of equations for the third time to the Academy. This time it was returned but with a negative report. The judges, who included Siméon-Denis Poisson, did not understand what Galois had written and (incorrectly) believed that it contained a significant error. They had been quite unable to accept Galois’s original ideas and revolutionary mathematical methods. The circumstances that led to Galois’s death in a duel in Paris are not altogether clear, but recent scholarship suggests that it was at his own insistence that the duel was staged and fought to look like a police ambush. In any case, anticipating his death the night before the duel, Galois hastily wrote a scientific last testament addressed to his friend Auguste Chevalier in which he summarized his work and included some new theorems and conjectures. Galois’s manuscripts, with annotations by Joseph Liouville, were published in 1846 in the Journal de Mathématiques Pures et Appliquées. But it was not until 1870, with the publication of Camille Jordan’s Traité des Substitutions, that group theory became a fully established part of mathematics. Britannica Web sites (1811-32). The work of a French mathematician, the prodigy Evariste Galois, was important in the development of modern algebra. His vital contributions to group theory, a part of higher algebra, showed the impossibility of trisecting the angle and squaring the circle, and answered many other long-standing questions. Article Contributors Article History Feedback Corrections? Updates? Help us improve this article!Contact our editors with your feedback.

Elementary school children's knowledge and intended behavior toward hearing conservation. The purposes of the study were to investigate children's knowledge about hearing conservation, the types of protective behaviors they would adopt in noise, the agreement between children's knowledge and intended behaviors in hearing protection, and reasons why they would not take any protective action in noise. A questionnaire was administered to 479 fourth and fifth graders in their school classrooms. Results indicated that children scored low (62.0%) on this hearing conservation questionnaire. They scored the highest in strategies of hearing protection (69.9%), followed by their knowledge in general hearing health (62.6%) and noise hazards (49.6%). Only 55% of children knew that hearing protective devices could protect them against noise. Approximately 28% of children did not intend to adopt any protective behavior in noise and the major reason for this was lack of knowledge. Children's knowledge and their noise protective behavior were correlated ( P < .05). However, possessing knowledge did not guarantee that children would adopt such behaviors when they were exposed to loud sounds. It is important to increase children's knowledge about hearing protection and hazardous noise as well as to encourage actual protective actions.

Q: Why I get this 0 as answer for (1/10) in Python? I am trying this in Python 2.7 Interpreter. >>> print 1/10 0 >>> print float(1/10) 0.0 >>> print 1/5 0 >>> a = 1/5 >>> a 0 >>> float(a) 0.0 >>> I want a value in floating point when I divide. Any idea on the logic behind this zero out put and please let me know how to do it correctly. A: The float conversion is happening after the integer division. Try this: >>> print(1.0/10) 0.1 A: In Python 2.x, by default / means integer (truncating) division when used with two integer arguments. In Python 3.x, / always means float division, and // always means integer division. You can get the Python 3.x behaviour in 2.x using: from __future__ import division See PEP 238 for a full discussion of the division operators. A: You need to do something like: print (1.0/10) When you do: print 1/10 It is evaluating it as 0 because it is treating it as an integer. And when you do: print float(1/10) You take the stuff inside the float() and evalute that first so in essence you get: print float(0) Which would be 0.0

Barriers to early detection of breast cancer among women in a Caribbean population. The purpose of this descriptive study was to identify and describe barriers to early detection of breast cancer, as well as current breast cancer screening behaviors and attitudes regarding the disease, among women aged 20 and older on the Caribbean island of Tobago. Tobago is the smaller of the two islands that make up the nation of Trinidad and Tobago. Between February and June 1996, 265 women fitting the age criteria completed a structured survey questionnaire. Women of African descent made up 89% of the respondents. In terms of age, 48% of those surveyed were between 20 and 39, 40% were between 40 and 59, and 12% were 60 or older. Barriers to early detection identified were a low level of breast self-examination, infrequent clinical breast examinations as part of regular care, unavailability of mammography services on Tobago, cost of screening, and difficulty of traveling to Trinidad for mammography. Furthermore, only a minority of the study participants had ever attended early detection or public awareness programs. The results were nearly the same for individuals with a family history of breast cancer, who would have a higher risk of occurrence of the disease. In addition, the majority of the respondents reported what can be considered a cultural barrier to early detection practices, a belief that no matter what they did, if they were to get breast cancer, they would get it. The authors recommend that mammography services be made available on Tobago. The authors also intend to use the findings to help develop an appropriate, culturally sensitive breast cancer awareness and early detection program for women on the island.

--- abstract: 'By networking multiple quantum devices, the Quantum Internet builds up a virtual quantum machine able to provide an exponential scale-up of the quantum computing power. Hence, the Quantum Internet is envisioned as a way to open a new distributed computing era. But the Quantum Internet is governed by the laws of quantum mechanics. Phenomena with no counterpart in classical networks, such as no-cloning, quantum measurement, entanglement and teleporting, impose very challenging constraints for the network design. Specifically, classical network functionalities, ranging from error-control mechanisms to overhead-control strategies, are based on the assumption that classical information can be safely read and copied. But this assumption does not hold in the Quantum Internet. As a consequence, the design of the Quantum Internet requires a major network-paradigm shift to harness the quantum mechanics specificities. The goal of this work is to shed light on the challenges and the open problems of the Quantum Internet design. To this aim, we first introduce some basic knowledge of quantum mechanics, needed to understand the differences between a classical and a quantum network. Then, we introduce quantum teleportation as the key strategy for transmitting quantum information without physically transferring the particle that stores the quantum information or violating the principles of quantum mechanics. Finally, we describe the key research challenges to design quantum communication networks.' author: - '\' - '\' bibliography: - 'quantum-04.bib' title: 'Quantum Internet: Networking Challenges in Distributed Quantum Computing' --- Quantum Internet, Quantum Network, Entanglement, Teleporting. Introduction {#sec:1} ============ Nowadays, the development of quantum computers is experiencing a major boost, since tech giants entered the quantum race. In November 2017 IBM built and tested a 50-qubits processor, in March 2018 Google announced a 72-qubits processor, and other big players, like Intel and Alibaba, are actively working on double-digit-qubits proof-of-concepts. Meanwhile, in April 2017 the European Commission launched a ten-years 1 €-billion flagship project to boost European quantum technologies research. And in June 2017 China successfully tested a 1200 km quantum link between satellite Micius and ground stations. Such a race in building quantum computers is not surprising, given their potential to completely change markets and industries - such as commerce, intelligence, military affairs [@CalCacBia-18; @PirBra-16; @Kim-08]. In fact, a quantum computer can tackle classes of problems that choke conventional machines, such as molecular and chemical reaction simulations, optimization in manufacturing and supply chains, financial modelling, machine learning and enhanced security. The building block of a quantum computer is the quantum bit (qubit), describing a discrete two-level quantum state as detailed in the next section. By oversimplifying, the computing power of a quantum computer scales exponentially with the number of qubits that can be embedded and interconnected within [@Bou17; @CalCacBia-18; @Kim-08]. The greater is the number of qubits, the harder is the problem that can be solved by a quantum computer. For instance, solving some fundamental chemistry problems is expected to require[^1] “hundreds of thousands or millions of interconnected qubits, in order to correct errors that arise from noise” [@Bou17]. Quantum technologies are still far away from this ambitious goal. In fact, so far, although the quantum chips storing the qubits are quite small, with dimensions comparable to classical chips, they usually require to be confined into specialized laboratories hosting the bulked equipment – such as large near absolute-zero cooling systems – necessary to preserve the coherence of the quantum states. And the challenges for controlling, interconnecting, and preserving the qubits get harder and harder as the number of qubits increases. Currently, the state-of-the-art of the quantum technologies limits this number to double digits (IBM 50-qubits and Google 72-qubits). Hence, very recently, the *Quantum Internet* has been proposed as the key strategy to significantly scale up the number of qubits [@CalCacBia-18; @Kim-08]. Specifically, by interconnecting multiple quantum devices via a *quantum network*, i.e., through a network able to share quantum information among remote nodes, and by adopting a distributed paradigm, the Quantum Internet builds up a virtual quantum machine constituted by a number of qubits that scales with the number of interconnected remote devices. This, in turn, implies an exponential speed-up of the quantum computing power [@CalCacBia-18]. From an engineering communication perspective, the design of the Quantum Internet constitutes a breakthrough. In fact, a quantum network is governed by the laws of quantum mechanics. Hence, phenomena with no counterpart in classical networks – such as no-cloning, quantum measurement, entanglement and teleporting – imposes terrific constraints for the network design. As instance, classical network functionalities, such as error-control mechanisms (e.g., ARQ) or overhead-control strategies (e.g., caching), are based on the assumption that classical information can be safely read and copied. But this assumption does not hold in a quantum network. As a consequence, the design of a quantum network requires a major paradigm shift to harness the key peculiarities of quantum information transmission, i.e., entanglement and teleportation. The goal of this paper is to shed light on the challenges in designing a quantum network. Specifically, we first introduce some basic knowledges of quantum mechanics – superposition principle, no-cloning theorem, quantum measurement postulate and entanglement – needed to understand the differences between a classical and a quantum network. Then, we focus on quantum teleportation as *the key strategy* for transmitting quantum information without either the physical transfer of the particle storing the quantum information or the violation of the quantum mechanics principles. Finally, we draw the key research challenges related to the design of a quantum network. Quantum Mechanics Background {#sec:2} ============================ In this section, we briefly review some quantum mechanics postulates and principles needed to understand the challenges behind the design of a quantum network. A quantum bit, or *qubit*, describes a discrete[^2] two-level quantum state, which can assume two (orthogonal) basis states: the zero (or ground state) and the one (or excited state), usually denoted[^3] as $\ket{0}$ and $\ket{1}$. As instance, if we represent the state of a photon with a qubit, the two basis states represent the horizontal and the vertical polarization of the photon, respectively. Superposition Principle {#sec:2.1} ----------------------- As widely known, a classical bit encodes one of two mutually exclusive states, being in only one state at any time. Conversely, a qubit can be in a superposition of the two basis states, being so simultaneously zero and one at a certain time. As instance, a photon with degrees of polarization is described by a superposed qubit, with an even *amount* of zero and one, being simultaneously horizontally and vertically polarized. Hence according to the superposition principle, $n$ qubits can simultaneously encode **all** the $2^n$ possible states at once. Differently, $n$ classical bits can encode only **one** of $2^n$ possible states at a certain time. Quantum Measurement {#sec:2.2} ------------------- According to one of the quantum mechanics postulates, whenever a measurement can have more than one outcome, as is the case for the two possible states of a qubit, after the measurement the original quantum state collapses in the measured state. Hence, the measurement alters irreversibly the original qubit state [@NieChu-11]. And, the result of such a measurement is probabilistic, since one obtains either the state zero or the state one, with a probability depending on the *amount* of zero and one in the original superposed quantum state. For instance, if the outcome of measuring a superposed qubit is the state zero, the qubit collapses into such a state and any further measurement will give zero as outcome, independently of the original *amount* of one in the superposed state. As a consequence, although a qubit can *store* more than one classical bit of information thanks to the superposition principle, by measuring a qubit, only one bit of information can be obtained. The measurement postulate has deep implications on the quantum network design as described in Sec. \[sec:4\]. In fact, we can not share quantum states among remote quantum devices by simply measuring the qubits and transmitting the measurement outcomes. Instead, we must share qubits among remote devices without measuring them by exploiting a fundamental quantum mechanics property: *the entanglement*. No-cloning theorem {#sec:2.3} ------------------ The no-cloning theorem states that an unknown qubit cannot be cloned. The no-cloning theorem is a direct consequence of the properties of the transformations allowed in quantum mechanics. Specifically, nature does not allow arbitrary transformations of a quantum system. Nature forces these transformations to be unitary. The linearity of unitary transformations alone implies the no-cloning theorem [@NieChu-11]. The *no-cloning theorem* has critical consequences from a communication engineering perspective, since classical communication functionalities are based on the assumption to be able to safely copy the information. This in turn deeply affects the quantum network design, as pointed out in Sec. \[sec:4\]. ![Quantum Entanglement: measuring one qubit of an EPR pair instantaneously changes the status of the second qubit, regardless of the distance dividing the two qubits. Specifically, the two particles forming the EPR pair are in superposed state, with an even *amount* of zero and one. Hence, by measuring one of the two qubits, we would obtain either zero or one with even probability. But once the qubit is measured, say particle A with outcome corresponding to state zero, the second qubit instantaneously collapses into state zero as well.[]{data-label="fig:09"}](./Fig-09.png){width=".8\columnwidth"} Entanglement {#sec:2.4} ------------ The deepest difference between classical and quantum mechanics lays in the concept of *quantum entanglement*, a sort of correlation with no counterpart in the classical word. In a nutshell, the entanglement is a special case of superposition of multiple qubits where the overall quantum state can not be described in terms (or as a tensor product) of the quantum states of the single qubits. To better understand the entanglement concept, let us consider Figure \[fig:09\], showing a couple of maximally entangled qubits, referred to as *EPR pair*[^4]. The couple of qubits forming the EPR pair are in a superposed state, with an even amount of zero and one. By measuring each of the two qubits independently, one obtains a random distribution of zero and one outcomes with equal probability. However, if the results of the two independent measurements are compared, one observes that every time that the measurement of qubit A yielded zero so did the measurement of qubit B, and the same happened with the outcome 1. Indeed, according to quantum mechanics as soon as one of the two qubits is measured the state of the other is *instantaneously* fixed. This quantum entanglement behavior induced Einstein and his colleagues to the so-called *EPR paradox*: the measurement of one qubit instantaneously changes the state of the second qubit, regardless of the distance dividing the two qubits. This seems involving information being transmitted faster than light, violating so the Relativity Theory. But the paradox is apparent, since entanglement does not allow to transmit information faster than light, as shown in Sec. \[sec:3.1\]. {width="1\columnwidth"} \[fig:11.a\] {width="1\columnwidth"} \[fig:11.b\] Quantum Communications {#sec:3} ====================== {width="1\columnwidth"} \[tab:01\] So far, quantum communications have been widely restricted to a synonymous of a specific application, e.g. *Quantum Key Distribution* (QKD) [@NieChu-11]. QKD is a cryptographic protocol enabling two parties to produce a shared random secret key by relying on the principles of quantum mechanics, either quantum measurement or entanglement. However, in a QKD system, quantum mechanics plays a role only during the creation of the encryption key: the encrypted information subsequently transmitted is entirely classical. Similarly, *superdense* coding [@NieChu-11] is a communication protocol enabling two parties to enhance the transmission of classical information through a quantum channel, i.e., to exchange two bits of classical information by exchanging a single qubit. Both QKD and superdense coding exploit quantum mechanics to convey classical information. Differently, as shown in Table \[tab:01\], the Quantum Internet relies on the ability to share quantum states among remote nodes. However, quantum mechanics restricts a qubit from being copied or safely measured. Hence, although a photon can encode a qubit and it can be directly transmitted to a remote node, e.g., via a fiber link, if the traveling photon is lost due to attenuation or corrupted by noise, the quantum information is definitely destroyed. This quantum information can not be recovered via a measuring process and/or a copy of the original information, due to the measuring postulate and the no-cloning theorem. As a consequence, the direct transmissions of qubits via photons is not feasible and quantum teleportation, described in the following, must be employed. Quantum Teleportation {#sec:3.1} --------------------- Quantum Teleportation [@BenBraCre-93] provides an invaluable strategy for transmitting qubits without either the physical transfer of the particle storing the qubit or the violation of the quantum mechanics principles. Indeed, with just local operations, referred to as Bell[^5]-State Measurement (BSM)[^6], and an EPR pair shared between source and destination, quantum teleportation allows one to “*transmit*” an unknown quantum state between two remotes quantum devices. Differently from what the name seems to suggest, quantum teleportation implies the destruction of both the original qubit (encoding the quantum information to be transmitted) and the entanglement-pair member at the source, as a consequence of a measurement operation. Then, the original qubit is *reconstructed* at the destination once the output of the BSM at the source – 2 classical bits – is received through a classical channel. The teleportation process of a single qubit is summarized in Figure \[fig:11\]. In a nutshell, it requires: i) the generation and the distribution of an EPR pair between the source and destination, ii) a classical communication channel for sending the two classical bits resulting from the BSM measurement. Hence, it is worth noting that the integration of classical and quantum resources is a crucial issue for quantum networks. Regarding the EPR pair, the measurement at the source side destroys the entanglement. Hence, if another qubit needs to be teleported, a new EPR pair must be created and distributed between the source and the destination. Quantum Network Design: Challenges and Open Problems {#sec:4} ==================================================== In this Section, we discuss the key research challenges and open problems related to the design of a quantum network, which harnesses quantum phenomena with no-counterpart in the classical reality, such as entanglement and teleportation, to share quantum states among remote quantum devices. Decoherence and Fidelity {#sec:4.0} ------------------------ Qubits are very fragile: any interaction of a qubit with the environment causes decoherence, i.e., a loss of information from the qubit to the environment as time passes. Clearly, a perfectly-isolated qubit preserves its quantum state indefinitely. However, isolation is hard to achieve in practice given the current maturity of quantum technologies. Furthermore, perfect isolation is not desirable, since computation and communication require to interact with the qubits, e.g., for reading/writing operations. Decoherence is generally quantified through *decoherence times*, whose values largely depend on the adopted technology for implementing qubits. For qubits realized with superconducting circuits[^7], the decoherence times exhibit order of magnitude within $10$-$100 \,\mu$s [@BenKel-15; @FilSheMag-15]. Although a gradual decrease of the decoherence times is expected with the progress of the quantum technologies, the design of a quantum network must carefully account for the time constraints imposed by the quantum decoherence. Decoherence is not the only source of errors. Errors practically arise with any operation on a quantum state due to imperfections and random fluctuations. Here, a fundamental figure of merit is the *quantum fidelity*. The fidelity is a measure of the distinguishability of two quantum states, taking values between 0 and 1. The larger is the imperfection of the physical implementation of an arbitrary quantum operation, the lower is the fidelity. Since teleportation consists of a sequence of operations on quantum states, the imperfection of such operations affects the fidelity of the “transmitted" qubit. From a communication engineering perspective, the joint modeling of the errors induced by the quantum operations together with those induced by entanglement generation/distribution (as described in the next subsection) is still an open problem. Furthermore, the no-cloning theorem prevents the adoption in quantum networks of classical error correction techniques, which depend on information cloning, to preserve quantum information against decoherence and imperfect operations. Recently, many quantum error correction techniques have been proposed as in [@Hanzo18]. However, further research is needed. In fact, quantum error correction techniques must handle not only bit-flip errors, but also phase-flip errors, as well as simultaneous bit- and phase-flip errors. Differently, in the classical domain, a single type of error, i.e., the bit-flip error, has to be considered. From the above, the counteraction of the errors induced by decoherence and imperfect quantum operations in a quantum network is a functionality embracing aspects that traditionally belong to the physical layer of the classical network stack. Entanglement Distribution {#sec:4.1} ------------------------- ![Entanglement Swapping. Two EPR pairs are generated and distributed: i) between a source (Alice) and an intermediate node (Quantum Repeater), and ii) between the intermediate node and a destination (Bob). By performing a BSM on the entangled particles at the Quantum Repeater, entanglement is eventually generated between Alice and Bob.[]{data-label="fig:12"}](./Fig_3_new_2.png){width="0.9\columnwidth"} According to the quantum teleportation protocol, as in classical communication networks, the “transmission” of quantum information is limited by the classical bit throughput, necessary to transmit the output of the BSM outcome. But, differently from classical networks, the “transmission” of quantum information is achievable only if an EPR pair can be distributed between remote nodes. Long-distance entanglement distribution – although deeply investigated by the physics community in the last twenty years – still constitutes a key issue due to the decay of the entanglement distribution rate as a function of the distance. There is a global consensus toward the use of photons as entanglement carriers, i.e., as candidates to generate entangled pairs among remote devices. Figure \[fig:12\] reports a possible strategy for entanglement distribution and we refer the reader to [@PirEisWee-15] for a survey. Entanglement distribution is achieved through Quantum Repeaters [@VanTou-13], i.e. devices that implement a physical process known as *entanglement swapping* [@VanTou-13]. In practice, two EPR pairs are generated with source (Alice in Figure \[fig:12\]) and destination (Bob in Figure \[fig:12\]) receiving one element of each pair while the other two are sent to an intermediate node (the Quantum Repeater in Figure \[fig:12\]). By performing a BSM on the two entangled particles at the intermediate node, entanglement is created between the elements at the remote nodes. The process does not require the elements of the two pairs to have the same nature, therefore it can be used to distribute entanglement between two remotely located atoms or ions or to enhance the useful distance for the distribution of entangled photons. Hence, instead of distributing the entanglement over a long link, the entanglement is distributed iteratively through smaller links. Despite these efforts, further research is needed. Specifically, from a communication engineering perspective, the entanglement distribution is a key functionality of a quantum network embracing aspects belonging to different layers of the classical network stack. More in detail: - *Physical Layer*. The entanglement is a perishable resource. Indeed, due to the inevitable interactions with the external environment, the entanglement among entangled parties is progressively lost over time. Entanglement losses deeply affect the quality of the teleportation process. For the reasons highlighted in the previous subsection, classical error correction techniques cannot be used to counteract these losses. Hence, robust entanglement distribution techniques based on quantum error correction are mandatory for the deployment of a quantum network. This still represents a very hard challenge in this field. - *Link Layer*: The no-broadcasting theorem, a corollary of the no-cloning theorem, prevents quantum information from being transmitted to more than a single destination. This is a fundamental difference with respect to classical networks, where broadcasting is widely exploited for implementing several layer-2 and -3 functionalities, such as medium access control and route discovery. As a consequence, the link layer must be carefully re-thought and re-designed to distribute the entanglement among multiple quantum devices. Furthermore, effective multiplexing techniques for quantum networks should be designed to allow multiple quantum devices to be connected to a single quantum channel (e.g. a fiber), and the access to the medium could be based on photon-frequency-division for the entanglement distribution. - *Network Layer*: The entanglement distribution determines the connectivity of a quantum network in term of capability to perform the teleporting among the involved quantum devices. Hence, novel quantum routing metrics are needed for effective entanglement-aware path selection. Furthermore, the teleportation process destroys the entanglement as a consequence of the BSM at the source side. Hence, if another qubit needs to be teleported, a new entangled pair needs to be created and distributed between the source and the destination. This constraint has no-counterpart in classical networks and it must be carefully accounted for an effective design of the network layer. Matter-Flying Interface {#sec:4.2} ----------------------- {width="85.00000%"} As mentioned before, there exists a general consensus about the adoption of photons as substrate for the so-called *flying* qubits, i.e., as entanglement carriers. The rationale for this choice lays in the advantages provided by photons for entanglement distribution: weak interaction with the environment (thus, reduced decoherence), easy control with standard optical components as well as high-speed low-loss transmissions. The aim of the flying qubits is to “transport” qubits out of the physical quantum devices through the network for conveying quantum information from the sender to the receiver. Hence, as shown in Figure \[fig:07\], a transducer [@Ham10] is needed to convert a *matter* qubit, i.e., a qubit for information processing/storing within a computing device, in a flying qubit, which creates entanglement among remote nodes of the network. Nowadays, there exist multiple technologies for realizing a matter qubit (quantum dots, trasmon, ion traps, etc) and each technology is characterized by different pros and cons, as surveyed in [@VanDev-16]. As a consequence, a matter-flying interface is required to face with this technology diversity. Moreover, from a communication engineering perspective, the interface should be compatible also with the peculiarities of the physical channels the flying qubits propagate through. In fact, there exist different physical channels for transmitting flying qubits, ranging from free-space optical channels (either ground or satellite free-space) to optical fiber channels. In the last ten years, the physics community has been quite active on investigating schemes and technologies enabling such an interface, with an heterogeneity of solutions [@Ham10; @StaRabSor-10] ranging from cavities for atomic qubits to optomechanical transducers for superconducting qubits. As a consequence, the communication engineering community should join these efforts by designing communication models that account both the technology diversity in fabricating qubits and the propagation diversity in characterizing the different physical channels. Deployment Challenges {#sec:4.3} --------------------- Quantum Internet is probably still a concept far from a real world implementation, but it is possible to outline a few deployment challenges for the near future. - The current technological limits to qubits and quantum processors physical realizations:\ At first, quantum computers will be available in few, highly specialized, data centers capable of providing the challenging equipment needed for quantum computers (ultra-high vacuum systems or ultra low temperature cryostats). Companies and users will be able to access quantum computing power as a service via cloud. Indeed, the quantum cloud market is estimated nearly half of the whole 10 billion quantum computing market by 2024 [@HSRC-18]. IBM is already allowing researchers around the world to practice quantum algorithm design through a classical cloud access to isolated 5-, 16- and 20-qubits quantum devices. - Existing technological limits to quantum communication and quantum interfaces:\ The first realizations of a Quantum Internet will be small clusters of quantum processors within a data center. Architectures will have to take into account the high cost of data buses (economically and in term of quantum fidelity) limiting both the size of the clusters and the use of connections for processing. - Hybrid architectures will probably be used for connections faring both the use of cryo-cables (expensive and necessarily limited in length) and of optical fibers or free space photonic links. - Given the fragile nature of quantum entanglement and the challenges posed by the sharing of quantum resources, a substantial amount of conceptual work will be needed in the development both of novel networking protocols and of quantum and classical algorithms. In conclusion, the Quantum Internet, though still in its infancy, is a very interesting new concept where a whole new set of novel ideas and tools at the border between quantum physics, computer and telecommunications engineering will be needed for the successful development of the field. [Angela Sara Cacciapuoti]{} (M’10-SM’16) is a Tenure-Track Assistant Professor at the University of Naples Federico II, Italy. Since July 2018, she held the national habilitation as “Full Professor” in Telecommunications Engineering and, since April 2017, she held the habilitation as “Associate Professor”. Her work appeared in first tier IEEE journals, and she received several awards. Currently, Angela Sara serves as Editor/Associate Editor for: IEEE Trans. on Communications, IEEE Communications Letters, and IEEE Access. In 2016 she was elevated to IEEE Senior Member and she has been an appointed member of the IEEE ComSoc YP Standing Committee. Since 2018, she is the Publicity Chair of the IEEE ComSoc WICE Standing Committee. [Marcello Caleffi]{} \[M’12, SM’16\] is a tenure-track assistant professor at the University of Naples Federico II. His work has appeared in several premier IEEE transactions and journals, and he has received multiple awards, including best strategy, most downloaded article, and most cited article awards. Currently, he serves as editor for IEEE Communications Magazine, IEEE Access, and IEEE Communications Letters. He serves as Distinguished Lecturer for the IEEE Computer Society and, since 2017, he serves as elected treasurer for IEEE ComSoc/VT Italy Chapter. [Francesco Tafuri]{} is a Full Professor at the University of Naples Federico II. His scientific interests fall in the field of Superconductivity. He is leading an experimental research group investigating macroscopic quantum phenomena and the Josephson effect in hybrid devices, new routes for the realization of superconducting qubits for hybrid architectures, mesoscopic and nanoscale systems, vortex physics. He is Associate Editor of Journal of Superconductivity, in the Advisory Board of Physica C, and Peer Reviewer of National and European projects. He is author of more than 160 publications on international journals and books, including 13 papers on Science/Nature/Nat.Mat./Nat.Comm./PRL. He has given $>$90 invited talks at international conferences and research institutes. Several manuscripts are the product of international and national collaborations, including IBM T.J. Watson Research Center, USA; Stanford University, USA; Chalmers University of Technology, Goteborg, Sweden; University of Cambridge, UK; Lawrence Berkeley Laboratory, USA; Columbia University, USA. [Francesco Saverio Cataliotti]{} activity is mainly experimental and is centered on atomic physics and atom-laser interactions. It can be divided in lines that, starting form nonlinear atom-laser interaction, develop into the manipulation of the atomic momentum state, laser-cooling and the realization of Bose-Einstein condensates. The optical manipulation of condensates has allowed him to investigate problems concerning superfluidity and macroscopic quantum states. More recently he has been concerned with coherent micro-manipulation techniques with atoms, molecules and micro-mechanical systems with perspectives ranging from quantum computation to single molecule manipulation. He is author of more than 100 scientific papers on international peer reviewed journals and contributions to conferences and workshops. Overall his works have received more than 2900 citations to date (researcherID K-4772-2015; orcid.org/0000-0003-4458-7977). [Stefano Gherardini]{} was born in 1989. He received his B.Sc. degree in Electronic and Telecommunications Engineering and his M.Sc. degree in Electrical and Automation Engineering from the University of Florence, Italy, in 2011 and 2014, respectively. He got his Ph.D. degree in Information Engineering, curriculum non-linear dynamics and complex systems, in 2018 with the Automatic Control and QDAB group at the University of Florence. Moreover, he is a visiting student at SISSA, the International School for Advanced Studies, in Trieste, Italy. His main research interests include binary measurements, networked moving-horizon estimation, Kuramoto models, disordered systems, stochastic quantum Zeno effect, open quantum systems theory, and quantum estimation and control. [Giuseppe Bianchi]{} has been a full professor of networking at the School of Engineering of the University of Roma “Tor Vergata” since 2007. His research activity, documented in more than 200 peer-reviewed international journal/conference papers having received to date more than 14,000 citations, focuses on wireless networks (his WLAN modeling work received the ACM SIGMO-BILE Test Of Time award in 2017), programmable network systems, and privacy and security monitoring in both the wireless and wired domains. [^1]: Qubit redundancy is mandatory due to the unavoidable presence of an external noisy environment. [^2]: At nano-scale, physical phenomena is described by discrete quantities, such as the energy levels of electrons in atoms or the horizontal/vertical polarization of a photon. This discrete nature contrasts with classical phenomena, which is commonly described by continuous quantities. [^3]: The *bra-ket* notation $\ket{\cdot}$, introduced by Dirac, is a standard notion for describing quantum states. In a nutshell, a ket $\ket{\cdot}$ represents a column vector, hence the standard basis $\ket{0},\ket{1}$ is equivalent to a couple of 2-dimensional orthonormal vectors. [^4]: A couple of maximally entangled qubits is called EPR pair in honor of an article written by Einstein, Podolsky, and Rosen in 1935. [^5]: So called from John S. Bell who, in 1964, proposed a physical test of the EPR paradox. [^6]: We emphasize that a solid strategy for a complete BSM with arbitrary photon states is not yet available. [^7]: Superconducting qubits are based on Josephson tunneling junctions, which allow for a flow of electrons with a practically zero resistance between two superconductor layers. Under certain conditions, current can flow in both the directions, i.e., in a superposition of two states. As the two states are very close energetically, low temperatures (20 mK) are needed to avoid thermal fluctuations [@BenKel-15].

100 Interesting facts 100 Interesting Facts The Benedictine Monks ran the first regular ferry from Birkenhead to Liverpool. The Monks would row over to the fishing village in Liverpool on market days, and offer the service to travellers. The service was granted a Royal Charter by Edward III in 1330. Edward III also granted the right to the Earl of Chester to run the ferry service from Seacombe to Birkenhead, establishing the Wallasey ferry. These two operations merged as Mersey Ferries in 1968. Rowing across the Mersey would take 90 minutes in calm weather; however much longer in rough conditions. In one day, the ferry had over 100 arrivals and departures on the Liverpool – Wirral service. The ferry used to run all night before the opening of the Mersey tunnels. The Mersey Ferries were the only way of getting across the river until the opening of the railway tunnel in 1886. In 1863 when the Channel fleet visited, 55,000 people took the ferry out to see the Navy ships. Before radar was installed in 1947, ferry captains had to rely on fog bells to give them an audible target to aim for The Mersey Ferries were the first in the world to install a radar system for safe navigation in fog, in 1947. In 1357, the charge for a foot passenger was a halfpenny. The opening of the first underwater railway in 1886 and the lines electrification in 1903 was predicted to be the end of the Mersey Ferries, but in 1919 the Seacombe ferry alone carried 22 million passengers. After the opening of the road tunnel in 1934, there was a huge decline in numbers using the Mersey Ferries. The opening of the second Mersey tunnel in 1971 brought the closure of the New Brighton ferry. The Birkenhead and Wallasey ferry operations merged into Mersey Ferries in 1968. The Mersey Ferries were under threat in 1977, when a bill was put before Parliament to discontinue ferry services. The bill failed to receive support and the Mersey Ferries survived. After a decline in commuter traffic using the ferries, in 1990 Merseytravel re-launched the ferries as a heritage and visitor attraction. Mersey Ferries won Visitor Attraction of the Year award in 1996, by the Merseyside Tourism Board. The Mersey Ferries took a technological leap in 1815 with the introduction of the first steam ferry, Etna. For the first time, the ferry service could operate a timetable. Soon after Etna, came Vesuvius, the second paddlesteamer. By 1840, all ten ferry services across the Mersey were using paddlesteamers. Crocus, the first coal-fired screw steamer was introduced in 1884, closely followed by Snowdrop the following year. The first diesel-powered ferry came in 1949 The last steam ferry on the Mersey was the Wallasey, which sailed for the last time at Whitsun in 1963. The template for the 20th century Mersey Ferries was the Wallasey ship Rose, in 1900. Cammell Laird built all five ferries that made up the Birkenhead post-war fleet; Upton, Hinderton, Thurstaston, Claughton and Bidston. Birkenhead and Wallasey followed different trends in naming their ferries; Birkenhead used local place names such as Woodchurch and Claughton, and Wallasey used flowers, such as Daffodil and Primrose. In 1906 Wallasey took ownership of Iris and Daffodil, which become their most famous ferries. In 1918, Daffodil and Iris were requisitioned by the Admiralty to take part in a raid on Zeebrugge on St Georges Day, 23 April. A commemorative service for the Zeebrugge raid is held every year on a Mersey ferry, on the Sunday nearest to St Goerge’s Day. Eleven Victoria Crosses were awarded for bravery in the Zeebrugge raid. In recognition of their courageous part in the raid, King George V granted Wallasey Corporation permission to add the ‘Royal’ prefix to both names. In November 1998, Overchurch went into dry dock in Manchester, where she was refurbished and given new engines. She later emerged as the Mersey Ferries new flagship, Royal Daffodil. Mountwood went into dry dock for refurbishment in 2001, and was renamed Royal Iris of the Mersey on her return in April 2002. The current Mersey Ferries fleet consists of two boats: Snowdrop, (previously Woodchurch) and Royal Iris of the Mersey (previously Mountwood). Up to 2013, the Royal Daffodil (previously Overchurch) had been the flagship of the fleet. The Royal Iris was sold in November 1991, and is now berthed in London, by the Thames Barrier. Since the name Royal Iris is still in use by the old ferry now berthed in London, we have Royal Iris of the Mersey. Each ferry holds 35 tons of fuel in its tanks and burns a ton in 12 hours: So, a Mersey ferry could travel across the Atlantic to New York without needing to refuel. Overchurch was the most recent Mersey ferry to be built, launched at Cammell Laird’s yard in 1962. The current ferries were built to carry 1,200 passengers; however they are now licensed to carry 650 for their daily service, and 396 for longer cruises. By 1989 the ferries were running at an annual loss of £2.5 million; a new strategy was then agreed, focusing on tourism and heritage rather than commuter transport. Between high and low tide, the river Mersey can rise and fall by 30ft. The ferries all have Fletner system rudders which make them much more manouverable. By using the vessel's twin screws, captains can move the vessels away by using one engine to push the vessel from the stage and the rudders and other engine to point it into the correct direction. A ferry across the Mersey from Seacombe, the narrowest crossing point of the river, is recorded in the Domesday Book of 1086 During the First World War the steamers Iris and Daffodil were taken out of service from Wallasey to be used as troop ships in the naval raid on Zeebrugge in Belgium The Manchester Ship Canal is the eighth-longest ship canal in the world. Mersey Ferries offer six hour cruises along the Manchester Ship canal with return bus transport. Famous people who have travelled on the Mersey Ferries – Queen Elizabeth and Prince Phillip, Princes Anne, Margaret Thatcher, King George V and Queen Mary, The Beatles, Gerry & the Pacemakers. The Manchester Ship Canal 60ft rise in height from Liverpool to Manchester is resolved by a series of 5 locks In 1847, the first floating landing stage opened at Liverpool. The landing stage rose and fell with the tide so that the boats could dock at any time. The first passenger ferry steamer to have a saloon was the Cheshire. It operated from Woodside in 1864. On 26 November 1878, the ferry Gem, a paddle steamer operated from Seacombe collided with the Bowfell, a wooden sailing ship at anchor on the River Mersey; the collision resulted in the death of five people. In 1894 trains were carrying 25,000 passengers per day and the ferries 44,000 per day. The ferry service at Tranmere closed in 1897. The pier and landing stage at Rock Ferry was built in 1899. In 1914 King George V and Queen Mary travelled on the ferry S.S. Daffodil from Wallasey to Liverpool. The ferry service from Seacombe lost two million passengers when the Queensway road tunnel opened in 1934, with people starting to use the tunnel rather than the ferry. The opening of the road tunnel also had an impact on the luggage boats which were introduced in 1879. Both ferry companies earned a substantial amount from luggage boats, which carried vehicles and goods across the river. The luggage boat services from both Woodside and Seacombe to Liverpool came to a close in the 1940s. The closure of Eastham in 1929 marked the last use of ferry paddle steamers on the river. In 1941, mines which had drifted into the River Mersey stopped ferry crossings In 1950, the ferries carried almost 30 million passengers. By 1970, the total number had fallen to just 7 million. In 1956, night boats across the river were cancelled and replaced by bus services through the tunnel. The future of the Mersey Ferries came under threat in June 1971, with the opening of the Kingsway road tunnel. The Mersey Ferries saw a further decline in passenger numbers; only 4,000-5,000 passengers a day using the service. 1984 was a considerably good year for the Mersey Ferries, and is seen as the beginning of the ferries rise from the slumps of the 1970s. A special ferry service was provided to Otterspool Promenade throughout the International Garden Festival. This service was usually operated by Overchurch. In August 1984, sailing ships from the Tall Ships Race visited the river. The first diesel ferry to enter service was the Royal Iris in 1951. The most famous Royal is the Royal Iris of 1951. She was by far the best loved of all the Mersey Ferries. Royal Iris was the first diesel powered vessel of the Wallasey fleet, with four diesel generators connected to two Metrovick marine propulsion units. Before being sold in the early nineties for use as a floating nightclub, The Royal Iris held hundreds of party cruises, and played host to popular bands such as Gerry & The Pacemakers, The Searchers, and The Beatles. The Royal Daffodil II was constructed by James Lamont and Co. at Greenock and entered service in 1957. The Royal Daffodil II is the only Mersey Ferry to ever have docking order telegraphs in the wings. Leasowe was sold to Greek owners in the mid 1970s; she has been heavily modified and is still cruising around the Greek islands. The former Royal Daffodil II was converted to a container ship. She sank in November 2007, 20 miles off the coast of Cape Andreas, in heavy seas. The Mountwood was used in the film "Ferry Cross The Mersey", The Mountwood was used for the video of the Gerry & The Pacemakers song, “Ferry across the Mersey” It was filmed on two separate journeys across to Liverpool from Birkenhead In her early years Mountwood was an unreliable ferry, breaking down three times whilst crossing the river and had to anchor. Her passengers had to be rescued by Woodchurch. In 1981, Woodchurch was withdrawn from operation for almost three years, due to cost cutting measures. Woodchurch returned to service in 1983, after main engine repairs and a full repaint. Her return freed up Overchurch to work the new Otterspool service, set up for the 1984 International Garden Festival. In 1996 the Overchurch was given a facelift, enclosing of the promenade deck shelter. The current fleet has served the river for nearly five decades. The Royal Iris of the Mersey and the Snowdrop celebrated their 50th birthday in 2009. The new Wartsilla engines fitted on the ferries are much more economical than the previous engines by Crossley Bros of Manchester. They produce much less emissions than the original propulsion units, making them a lot greener. Each of the Mersey Ferries in the current fleet carries two Kockums Super Tyfon TA 100/150 type fog horns. These are the original horns that were fitted when the ferries were first built. Both Royal Iris of the Mersey and Snowdrop have an E-flat tone, and Royal Daffodil is in F Sharp. The ferries are known for their ability to operate in very heavy seas. When berthing the vessel, the captain uses a combination of rudder positions and engine movements. Although the ferries can cope in heavy seas, berthing the vessel can be extremely hazardous. This can result in the service being suspended. The ferries all have Fletner system rudders which make them much easier to manoeuver. The ferries played a big part in Liverpool's European Capital of Culture 2008 celebrations, where they carried a record numbers of passengers On Monday 21st July 2008, the Royal Daffodil operated a special cruise to witness the parade of sail and departure of the Tall Ships. The first record of a service from New Ferry to Liverpool was in 1774. The first steamship to operate on the Mersey was the Elizabeth, a wooden paddle steamer, which was introduced in 1815 to operate between Liverpool and Runcorn. The Woodside, North Birkenhead and Liverpool Steam Ferry Company was formed in 1835. The Cheshire, the first passenger ferry steamer to have a saloon, operated from Woodside in 1864. The iron pier at Eastham was built in 1874. The pier and landing stage at Rock Ferry was built in 1899. Wallasey Ferries had a funnel livery of white and black and Birkenhead red and black up until the Second World War, when this was changed to orange when Mountwood and Woodchurch were introduced. Due to financial losses incurred from a reduction in patronage, Birkenhead Corporation gradually closed its southern terminals; New Ferry on 22nd September 1927, Eastham in 1929 and Rock Ferry on 30 June 1939. In 1941, The Oxton and Bebington vessels were fitted with cranes to enable them to unload United States aircraft from mid-river and deliver them to the Liverpool landing stage. The Birkenhead boats Claughton, Bidston, Thurstaston and Upton were seen as the fastest ferries on the river.

Capacitors of the type utilized with ballasts in AC applications are commonly made from a metallized film, such as polypropylene. To form the capacitor, the film is typically wound into one or more roll-shaped capacitor sections which are housed within a sealed metal container. Such containers often include a mechanical interrupter which disconnects the capacitor from external terminals, in the event that there is excessive pressure or temperature within the container. In addition, the capacitor is impregnated with oil and the entire container is filled with oil, in order to seal out air and moisture. This avoids the possibility that air or moisture inside the capacitor might ionize and produce electrical arcing. The problem with oil-filled capacitors is that the oils used are considered to present an environmental and health hazard. Accordingly, great care must be taken to avoid leakage of the oil from the capacitor casing while the capacitor is installed. In addition, disposal of capacitors that are no longer usable requires that they be subjected to special toxic waste treatment. In some instances, metallized capacitor sections for fluorescent lighting applications have been dipped in wax. However, wax exhibits a relatively high degree of shrinkage and has a tendency to develop cracks. In high voltage applications, air and/or moisture is likely to enter these cracks and results in early failure of the capacitor, owing to electrical arcing. Also, wax tends to be a relatively expensive material for use in capacitors. Broadly, it is an object of the present invention to avoid the shortcomings of existing capacitors used for a AC applications. It is specifically contemplated that the invention provide a capacitor used for relatively high voltage AC applications which does not require a liquid impregnant in order to seal the capacitor against the entry of air and moisture. It is another object of the present invention to provide a metal encased capacitor which does not require a liquid fill material, yet is resistive to failures due to the leakage of air and/or moisture thereinto, and which will continue to function reliably in relatively high voltage applications for extended periods of time. It is also an object of the present invention to provide a capacitor for AC applications which is reliable, safe and convenient in use, yet relatively inexpensive in construction. It is also an object of the present invention to provide a method for manufacturing a capacitor useful in AC applications which meets one or more of the foregoing objects. In accordance with a preferred embodiment of the present invention, a capacitor for an AC application is manufactured from zinc alloy metallized polypropylene sheet material. The sheet material is wound into one or more sections having a roll-shape, which sections are placed within a metal container and provided with electrical terminals for connecting to the exterior of the container. The metal container is then filled with a potting compound comprising 80%-90% blown asphalt and 10%-20% wax additive. When the potting compound sets, the capacitor sections are sealed therein so as to be protected from air, moisture and thermal and electrical shock. The container is then covered with a top which includes a pressure interrupter, and it is sealed. This capacitor exhibits improved resistance to deterioration due to environmental factors over long periods of time, and it permits operation with higher voltages than was previously possible.

When grocery shopping, everyone has a natural habit of picking and choosing the most aesthetically pleasing fruits and vegetables. However, sometimes the less attractive produce is equally as healthy and delicious. Supermarket food that is not purchased, often for superficial reasons, is thrown out, which creates food waste. To help reduce this problem, a successful Canadian supermarket retailer is selling its imperfect produce at lesser prices. Fruits and vegetables that are misshapen or considered “ugly,” but are otherwise good, are being sold on sale at Loblaws, Canada’s largest supermarket chain. This clever campaign is an intelligent way to raise awareness of the dangers of excessive food waste. Discarded food fills up landfills and decomposes causing the emission of greenhouse gases like methane, which is toxic to our environment. The campaign is also a way to help consumers save money on produce that is equal in nutrition to more “attractive” produce.

A Brillouin analysis sensor is a fibre optic distributed sensor that can measure changes in the local strain and/or temperature conditions of an optical sensing fibre through analysis of the Brillouin frequency of the sensing fibre at any point. The sensing fibre is usually attached to or embedded in a host material and is used to measure the strain and/or temperature conditions of the host. Both strain and temperature cause a shift in the Brillouin frequency. Position is determined by the round-trip transit time of the optical signal in the fibre. Reference herein to a Brillouin analysis sensor system shall mean a Brillouin analysis sensor plus an optical sensing fibre. A Brillouin analysis sensor system operates on the principle of Brillouin amplification. In a Brillouin amplifier, a signal (or Stokes) light wave propagating in one direction experiences optical gain if its frequency falls within the Brillouin gain profile of the amplifier pump wave propagating in the opposite direction. The Stokes wave experiences gain (with a corresponding loss from the pump) in accordance with the Brillouin gain profile. The power in either wave may be monitored to determine the Brillouin gain profile; if the Stokes wave is chosen to be monitored then the system is said to operate in Brillouin gain mode (as any interaction causes an increase in the Stokes wave). Otherwise, it is said to operate in Brillouin loss mode. Two lasers may be used to generate the two lightwaves, in which case the frequency difference between them is maintained and controlled by a control loop (often a phase-locked loop). Alternatively, a single laser source may be used and a second lightwave generated from it by optical frequency modulation (using an electronic signal of the desired frequency). One type of Brillouin analysis sensor is a Brillouin Optical Time Domain Analysis (“BOTDA”) system. To obtain spatial information in a BOTDA system, one of the two light waves is pulsed by an electro-optic modulator (EOM) driven by an electronic pulse generator, while the intensity of the other light wave is monitored as a function of time elapsed (i.e. in the time domain) since the generation of the pulse. In this way position, z, is determined from the time-of-flight, t, as z=ct/2n, where c is the speed of light and n is the group index of refraction of the fibre. The intensity of the interaction (either gain or loss depending on which wave is monitored) is a function of the relative frequency difference between the two lasers. It is also a function of the gain profile of the fibre at the point of interaction, which moves (with the pulse) along the fibre. The frequency difference is swept over some range of frequencies of interest. Since frequency characteristics of the Brillouin profile are modified by the strain and temperature conditions of the sensing fibre, the local strain and temperature conditions can be inferred from a measurement of the Brillouin gain or loss profile. This profile is obtained by considering the Brillouin interaction at a given point in space as a function of frequency difference. A second type of Brillouin analysis sensor is the Brillouin Optical Frequency Domain Analysis sensor (“BOFDA”). In frequency domain analysis, a sine-wave amplitude modulated light wave is used in place of the pulse modulated wave used in BOTDA. Measurements are carried out in the frequency domain by sweeping this modulation frequency and measuring the interaction with the cw wave. The frequency domain data can be mathematically transformed to the time domain to obtain the positional information which is the equivalent of the signal that would be obtained by BOTDA. A third type of Brillouin analysis sensor is the Brillouin Optical Correlation Domain Analysis sensor (“BOCDA”). In correlation domain analysis both light waves are modulated and positional information is extracted from the correlation of the two signals. Correlation based sensors can have very high resolutions, although at the cost of limited overall sensing length. Traditionally bright pulses have been used to interrogate the sensing fibre with BOTDA systems, however the use of dark pulses has recently been shown to have several advantages in terms of resolution, accuracy and speed. BOTDA (and all other Brillouin analysis distributed sensors) require that the state of polarization (“SOP”) be matched (parallel) between the two light waves (pump and Stokes) at all locations in the sensing fibre. If the SOP is not matched, then the intensity of the Brillouin signal will fade at locations where the SOP is mismatched, resulting in amplitude noise and possibly locations in the sensor where the two SOPs are orthogonal and no signal can be obtained. In prior art BOTDA sensor systems, the SOP can be kept aligned but only if special polarization-maintaining (“PM”) fibres are used as the sensing element. Such fibres have higher attenuation, are much more expensive and are more difficult to work with compared to standard single mode fibres (“SMF”). If non-PM fibre (i.e. SMF) is used, then the state of polarization of the two lightwaves will vary with position (and possibly with time) along the fibre length. To deal with this situation, a prior art BOTDA system will often take multiple measurements on SMF, altering the lightwaves' relative polarization states between measurements, and using an average result. It has been shown that using averaged data from switching one of the waves between two orthogonal polarization states can reduce polarization fading, however the averaged signal will only be one half of the maximum Brillouin signal. Alternatively, a set of measurements taken with three particular SOPs can improve this average to two thirds of the maximum, but requires at least triple the time to take data. Moreover, the devices typically used to alter the SOP in a BOTDA system have switching times which are slow enough to significantly add to the overall acquisition time of the sensor. Also, because the multiple measurements are made at different times, it is possible for changes to the conditions of the sensor (e.g. mechanical vibrations or motions) between measurements to alter the SOP so that the multiple measurements no longer average orthogonal SOPs, thus reducing the efficacy of polarization averaging methods. BOTDA distributed sensor systems are therefore subject to polarization fading effects which require expensive PM fibre or polarization averaging schemes.

Certain memories used in digital systems can experience data loss if too many physically-proximate data bits in the memory have the same bit value and surround a bit of opposite value. That is, if too many bits in physical proximity are zeros, data loss can occur for a one bit stored near the zeros. Similarly, if too many bits in physical proximity are ones, data loss can occur for a zero bit stored near the ones. To mitigate data loss in these types of memories, memory controllers may implement data randomization, also referred to as “data scrambling.” The data scrambling mechanisms use reproducible pseudo-random modifications to change data written to memory, and to change the data back to the original data when read from the memory. By changing the data, the mix of ones and zeros is changed in an attempt to reduce the number of like-bits in physical proximity.

Applying powders in a melt state onto surfaces--in particular metal surfaces--represents a well known method. So e.g. the HU-PS 159 923 discloses a method of this type for applying and preparing a coating. In course of said method the flame required for melting is produced by using a mixture of acetylene and oxygen or acetylene and air. Powder--e.g. aluminium powder--is blown by means of pressurized air along a plane between two slabshaped flames tending to each other and melt inbetween, thereafter the melt powder is applied onto the surface to be coated by using pressurized air. The disadvantage of said method and apparatus, respectively, lies in that acetylene is most expensive, energy exploitation is uneconomical, in addition, quality of coating obtained by the method and apparatus does not meet requirements. A process and an apparatus are also known which use PB-gas or natural gas as combustible gas. The disadvantage of this solution lies in that in order to prevent lateral deflection of the powder jet discharged in a slab-shape from the spraying head and thus to eliminate powder losses, slab-shaped cold air-jets are directed onto the powder jet. In order to prevent powder losses to the desired extent, a considerable quantity of cold air is required decreasing considerably the temperature of the flame needed for melting the powder, resulting in a lower efficiency of powder melting. As a consequence of the significant quantity of cold air energy required for melting increases considerably, stability of the melting flame decreases. Flame stability is further deteriorated by the fact that structural design of the spray head of known apparatuses does not exclude arrival of the air streams coming from different lateral directions and of uncontrollable intensity at the melting flame.

Hey I Know This Probably Sounds Like A Stupid Question We Are Learning About Pupillary Reactions Perla How Will The Pupils React To Light If Someone Is Blind In One Eye????? Thanks They should still both be equal in size. "Shining a light in one eye of a normal subject causes both pupils to constrict equally. The pupillary reaction in the illuminated eye is called the direct response, and the reaction in the other eye is the consensual response. Because of the hemidecussation of afferent pupillomotor fibers at the chiasm, and because there is another hemidecussation of the pupillomotor fibres in the brain stem, the direct and consensual pupillary responses are equal. If one eye is blind, all input to the pupillary centers in the brain stem comes from the other eye, but the double hemidecussation ensures equal pupillary innervation and there will be not anisocoria Shining a dim light in one eye of a normal subject will cause both puils to constrict. A brighter light will cause more constriction. If, after shining a light in one eye, the light is quickly switched to the other eye, the response will be an initial constriction of both pupils followed by an equivalent re-dilation." From "Neuro-ophthalmology. Basic and Clinical Science Course 1996-1997" American Academy of Ophthalmology. and "Pathophysiology: Pupil size is governed by the tone of the pupillary sphincter (parasympathetic) and the pupillary dilator muscles (sympathetic) in response to ambient light, adrenergic tone, and local pharmacologic or pathophysiologic conditions. Approximately 50% of the fibers of the optic nerve decussate in the optic chiasm, and the input to each of the parasympathetic nuclei in the brain stem remains equal. Under normal conditions, the pupils remain equal at all times in all levels of light. Therefore, anisocoria never is caused by retinal pathology. Even blindness in one eye does not cause anisocoria." This is an area where I feel WEAK. And now realizing my co- workers don't so much. Thank you so much for the info! I had 2 neuro pts last pm and in the ER and one was a 16yr old who had been assaulted. His pupils started out equal but became unequal. My second pt was legally blind and the RN going off said "I didn't do pupil checks, he's legally blind,they don't react." I thought he knew what he was talking about since he has been an ED RN longer than me. Guess not.

You'd think that penguins would love the meaty taste of fish, but it turns out that they may not be able to taste their food much at all. A new genetic study out of the University of Michigan finds that penguins appear to have long ago lost the ability to taste sweet and bitter flavors, as well as the savory, meaty flavor known as umami. Together, sweet, bitter, and umami make up three of the five basic tastes. The other two, sour and salty, may still be present in penguins. The findings are being published today in Current Biology. The fish may have been too cold for penguins to taste So what led to penguins' believed loss of these tastes? "This is the most difficult question, especially because penguins eat fish and fish have the umami taste, so you would have predicted that the umami taste would have been useful to penguins," George Zhang, the paper's corresponding author, tells The Verge. Zhang's hypothesis — and he notes "it is still a hypothesis" — is that it's too cold down in Antarctica for these taste receptors to matter. Prior research has shown that the tongue's receptor channels, which are what react to tastes, function poorly at lower temperatures when it comes to detecting sweet, bitter, and umami. The Michigan researchers speculate that the environment may have been so cold that this receptor was "effectively non-functional" in penguins' ancestors. "Those three tastes … would not be useful any more because the channel is not functioning," Zhang says. "So gradually mutations would accumulate in those genes and eventually they would become lost." The study is based entirely on genetic findings, so another group of researchers will have to actually run tests with penguins to confirm that they can't taste these foods. But the study's authors are quite confident that future work will show penguins can't taste sweet, bitter, and umami flavors based on the genetic findings. "It's very clear," Zhang says. If the genes aren't there, he says, then the animals don't have those tastes. The genes for tasting sweet, bitter, and umami were broken or missing "The neighboring genes were all there, it's just that those specific taste genes were missing," Zhang says. "So we know it's not because the quality of the genome sequence." Zhang's research group was able to locate the genes for tasting bitter and umami flavors in around 20 other birds, but in the five penguin genomes that they looked at, including those from Adélie and emperor penguins, the genes were either broken or missing. Birds as a group are known to have lost their ability to detect sweet tastes, so that gene was absent in every sample that was tested. The absence of bitter and umami tasting genes in the five penguins leads the researches to believe that penguins as a group likely lost this taste in their common ancestor, while the ability to detect sweet tastes appears to have been lost much farther back. Because penguins seem to have originated in the Antarctic, the researchers believe that penguins that live elsewhere should also be unable to detect these tastes, too. (Anton_Ivanov / Shutterstock) Though the researchers believe that penguins should still be able to detect sour and salty tastes, David Yarmolinsky, a taste researcher at Columbia who was not involved with the study, says that's still up in the air. "Really this paper can only make a statement about those three taste qualities," he says. That's because while the absence of a gene typically says something, its presence doesn't always mean the same thing. "They did find a channel that's necessary for salt taste … but that doesn't mean a whole lot because that channel is required for your kidneys to work," he says. "If they didn't have that, they'd be dead penguins whether or not they can taste. So just because they have it in their genome doesn't mean they have it on their tongue." Yarmolinsky also says that the gene associated with the ability to detect sour tastes has complications that prevent that tasting ability from being a certainty, too. "You really have to take the penguin and see how it reacts." "You really have to take the penguin and see how it reacts when you give it something sour to eat," he says. "That's really the test to see what they can or cannot taste." Earlier research suggests that what taste buds penguins do have are fairly limited, if they have them at all. That makes sense, since the main function of penguins' tongues appears to be allowing them to capture prey; penguins also swallow their food whole. "Given the way their mouths are oriented, and how they hunt, capture, and consume their prey, it’s not surprising that penguins may have limited taste perception," Michael Polito, an ecologist with Louisiana State University who was not involved with the study, writes in an email to The Verge. "Other senses, especially vision and possibly even smell, may be more important when it comes to how penguins find and identify their favorite foods in the ocean." Whales and dolphins are also believed to have lost most of their taste, with only the salt channel remaining. Yarmolinsky notes that researchers have been looking at taste loss and its consequences in other animals lately. This study, he says, "really fits into the story that's coming out the last few years that sensor receptors are adaptive to the lifestyles of animal diet." It's also believed that taste receptors may play other roles inside the body, so a penguin's inability to detect sweet, bitter, and umami flavors could have deeper implications than just what's on its tongue.

Methylmercury Biogeochemistry in Freshwater Ecosystems: A Review Focusing on DOM and Photodemethylation. Mercury contamination is a growing concern for freshwater food webs in ecosystems without point sources of mercury. Methylmercury (MeHg) is of particular concern, as this is the form of mercury that crosses the blood-brain barrier and is neurotoxic to organisms. Wetlands and benthic sediments have high organic content and low oxygen availability. Anaerobic bacteria drive the metabolic function in these ecosystems and subsequently can methylate mercury. The bioavailability of MeHg is controlled by physicochemical characteristics such as pH, binding affinities, and dissolved organic matter (DOM). Similarly, photodemethylation is influenced by similar characteristics and thereby the two processes should be studied in tandem. The degradation of MeHg through photochemistry is an effective destruction mechanism in freshwater lakes. This review will highlight the uncertainties and known effects of DOM on subsequent photoreactions that lead to the occurrence of mercury photodemethylation and reduction in mercury bioavailability in freshwater ecosystems.

WCR This Week Northern residential schools tried to save wilderness life This archival photo shows a birthday party being held at Holy Angels Indian Residential School in Fort Chipewyan, Alta. March 3, 2014 GLEN ARGAN WESTERN CATHOLIC REPORTER Catholic residential schools in northern Canada strove to preserve the traditional wilderness way of life, a leading Catholic educator wrote in a 1992 paper. "The schooling provided in the North was based on the assumption that most Métis and Indian children would take up a life of hunting, trapping and fishing, and that the three Rs, while not essential to such a mode of living, would nonetheless give them certain advantages," wrote Dr. Robert Carney, an education professor at the University of Alberta. Carney had served as a principal and superintendent of school programs in the Northwest Territories before becoming executive director of the Alberta Catholic School Trustees' Association and then moving on to the U of A. One of his four children is Mark Carney, governor of the Bank of England. Dr. Carney died in 2009. His paper, "Residential Schooling at Fort Chipewyan and Fort Resolution: 1874-1974," maintained that the schools saw the wilderness as being a continuing source of native subsistence and general well-being. "Accounts by former students and others indicate that the schools' goals were in keeping with the circumstances of the time and the students' own preferences," Carney wrote. The schools encouraged the relatively small number of students who attended them to adhere to their parents' Christian beliefs and follow their livelihoods as hunters and trappers, he said. In 1939, a maximum of 30 per cent of the eligible children attended the schools. Those objectives, he continued, were "severely compromised" by government programs in the 1950s and 1960s which encouraged native students to abandon life on the land and to seek wage employment elsewhere. NATIVE LANGUAGES Aboriginal languages were used to teach catechism in the schools at least through the 1930s, and students were encouraged to speak their native tongues in the schoolyard and other non-classroom contexts, he said. "Until this time, the religious, child care and basic skill programs at Chipewyan and Resolution were largely compatible with the tenets of the native-wilderness equation." Carney said classroom teaching by the religious sisters was seen as fair and effective. School inspectors and other visitors "invariably praised the sisters' teaching skills and the testimony of most former students is similarly appreciative." Nevertheless, the French-speaking sisters faced great obstacles. Many did not speak English – the language of instruction – before coming north, and there was no encouragement for them to learn the native languages. Classroom discipline, he said, was strict. Bishop Gabriel Breynat of the Vicariate of Mackenzie from 1902 to 1943 promoted the use of aboriginal languages in the schools and advocated practical programs involving bush skills. His pleas "fell on deaf ears" among government officials. WHITE SETTLERS Carney said the situation changed when white settlers successfully pushed for integrated, non-sectarian schools to be established in Chipewyan and Resolution by the late 1930s. "Schooling arrangements based on the native-wilderness equation were rejected out-of-hand by these groups." John Milloy, in his best-selling history of the residential schools, A National Crime, details federal government resistance to the native-wilderness approach to education. Milloy quotes another historian who maintained that until after the Second World War the government had no interest in eliminating the hunting and trapping lifestyles of northern natives. However, in the mid-1950s, the government began to speak in terms of modernizing, rather than eliminating, northern aboriginal culture. The goal was "to help them become better Indians and Eskimos," according to one document. Despite the rhetoric, Milloy said, "assimilation was the norm. The rhetoric of cultural sensitivity and preservation was not in the end matched by the reality of the system itself." CHANGING TIMES Nor could the education system live up to that ideal, he continued, given the government's policy of promoting northern economic development and establishing permanent communities of people who had previously lived on the land. He also quoted a memo by E.A. Coté, deputy minister of northern affairs, who pointed to the non-aboriginal future for residents of the North: "As attractive and as quaint as any aspect of culture may be, this does not justify an effort of freezing it into a state of perpetuity. "It would indeed constitute an anachronism in this age and in this setting if we sought to maintain a stone-age culture, for example, among people who find themselves impinged on every side by wave upon wave of modern technology." By this point, the native-wilderness equation was doomed even if, as Carney notes, Mr. Justice Thomas Berger temporarily reignited the dream in the mid-1970s with his famous report which led to the scuttling of the Mackenzie Valley Pipeline.

Why does an error message appear and Windows® fails to start? This symptom occurs when the system files (files necessary for starting of Windows) cannot be found. The following are some causes of this symptom. Note: User may wish to back up the data, when Windows restart to prevent any loss of data in future. 1. There is a floppy disk in the floppy disk drive.Remove the floppy disk and restart Windows. 2. StickingIf the VAIO computer is brought from a colder environment to a warmer one, a small amount of moisture in the air may condense on the head of the hard disk drive. This results in "sticking". If that happens, wait for a while. The condensed moisture will vaporize with the natural rise of temperature of the hard disk. Once the condensed mositure is vaporized, the system will start working normally. 3. The hard disk is not active or cannot be recognized. If FDISK was executed just before this symptom started occuring, the hard disk drive may not be activated. Execute FDISK again and select 4 Display partition information to check the status of the hard disk drive. If A (which means "Active") does not appear under the status, the settings are not done properly. Select Set active partition to activate C:1. Restart the computer. 4. System files do not exist in the hard disk drive.a. User can start Windows from the Recovery CD-ROM by inserting the Windows boot disk in the floppy disk drive. Execute the FDISK / MBR command to check if the system can recognize the hard disk drive. b. If step A does not work and user has the Windows boot disk, start Windows with the disk. Execute the SYS command (A:\>SYS C:[Enter]) to transfer the system files to the hard disk drive. c. Execute Recovery to install the system files to the hard disk drive.

[Identification of orthopoxvirus species using oligonucleotide micro-chips]. A method for describing the Orthopoxviruses that are pathogenic both to man and animals is described in the article. The method is based on hybridization of a fluorescently labelled amplified DNA sample with oligonucleotides, which were immobilized in a microchip. Species-specific regions within the crmB gene encoding a viral analogue of the tumor necrosis factor receptor, i.e. an important gene determining the pathogenicity of the mentioned Orthopoxviruses type, were used as a target for identification. The identification procedure takes around 6 hours and does not demand any costly equipment (a portable fluorescent microscope can be used).

Lindisfarne National Nature Reserve Lindisfarne National Nature Reserve is a UK national nature reserve. It was founded to help safeguard the internationally important wintering bird populations, and six internationally important species of wildfowl and wading birds winter here. For the pale-bellied brent geese from Svalbard, this is their only regular wintering place in all of the United Kingdom. Pinkfooted and greylag geese, wigeons, grey plovers and bar-tailed godwits are the other visitors. Ramsar Site Lindisfarne National Nature Reserve is a Ramsar site, and as a result is a wetland of international significance. Habitats Lindisfarne National Nature Reserve protects a stretch of coastline, including the dunes of Lindisfarne (Holy Island). Lindisfarne NNR has international recognition and covers a large and varied mosaic of internationally recognised and important coastal habitats. These include intertidal mudflats, rocky shore, sand dunes and salt marshes. The dunes of the Lindisfarne National Nature Reserve support many plants. Early forget-me-nots and marram grass are among those that hug the ground and need little water. Many species of insects, moths, and butterflies appreciate this environment, including the 'woolly bears' (tiger moth caterpillars). Dark green fritillary, and grayling butterflies bask in the sun during July and August in the Lindisfarne NNR. Big brown-lipped snails can also be found. During early summer, the purple northern marsh orchid flourishes along with its pinker relation, the early marsh orchid. In July the marsh helleborines flower by the thousands and form spectacular white carpets. Recently a unique orchid, the Lindisfarne helleborine, has been discovered on the island. The plants were formerly identified as dune helleborine (Epipactis dunensis), but DNA analysis carried out in 2003 revealed them to be genetically distinct, and the new species was given the scientific name Epipactis sancta. Birds Bird species for which the Lindisfarne National Nature Reserve is important include pale-bellied brent goose, wigeon, teal, pintail, merlin, dunlin, bar-tailed godwit and many others. The situation on the east coast also makes it a good place for observing migrating birds arriving from the east, including large numbers of redwing and fieldfare, and also scarcer Siberian birds including regular annual yellow-browed warblers. Rare species such as Radde's warbler, dusky warbler and red-flanked bluetail have all occurred on Lindisfarne. , 330 species have been recorded on the isle of Lindisfarne and the adjacent Lindisfarne National Nature Reserve. The large number and variety of birds present makes the area a very popular region with bird watchers, particularly in the autumn and winter. Common throughout Lindisfarne and the Farne Islands are the common eider duck. Saint Cuthbert was reputed to be very fond of them (as pets, not food) during his periods on the Farnes as a hermit. In 676 he promulgated a law protecting the birds, reputed to be the first bird protection legislation. Locally the birds are known as "Cuddy's ducks" (occasionally "Cuddy's hens") after the familiar form of "Cuthbert". Seals Grey seals are commonly found in the waters and on the rocks around the isle of Lindisfarne, and common seals are also occasionally seen. Location Lindisfarne National Nature Reserve is located on the North Northumberland coastline. It is located about south of Berwick-upon-Tweed. Lindisfarne is sign-posted from the A1 highway south of Berwick-upon-Tweed. There are car parks available on the isle of Lindisfarne (Holy Island), at Budle Bay, and at Beal, on the mainland of England. The causeway and safety Access to the island is by a tidal causeway. Visitors to the island must check tide times and weather carefully, and seek local advice if in doubt. The road is generally open from about 3 hours after high tide until 2 hours before the next high tide, but the period of closure may be extended during stormy weather. Walkers using the causeway over the mud flats are advised that the safe time is shorter. Despite these warnings, about one vehicle each month is stranded on the causeway. Tide tables giving the safe crossing periods are published by Northumberland County council. Volunteering opportunities The Lindisfarne National Nature Reserve offers a range of volunteering opportunities for youths, seniors, and families. These volunteer opportunities include: Species Recording, Habitat and Estate Management, Site Wardening, and acting as guides during guided walks. References Category:National nature reserves in England Category:Ramsar sites in England Category:Geography of Northumberland Category:Sites of Special Scientific Interest in Northumberland Category:Nature reserves in Northumberland Category:Lindisfarne

Kids learn stroke signs in class through imitation In this Feb. 25, 2014 photo, third-graders, from left, Hunter Thomas, Matthew Velez, Sebastian Mendez, Jayden Gonzalez and Elijah Farias examine a plastic model of a brain at Montefiore Hospital in the Bronx borough of New York. The hospital teaches children to recognize stroke and get victims to a hospital quickly. (AP Photo/Jim Fitzgerald) NEW YORK (AP) — Andrea Esteban tried to smile with half her face, crossing her eyes in the process, and her third-grade classmates giggled. Matthew Velez struggled to speak, “Luh, luh, uh, gronk,” and the kids erupted in laughter. But the funny faces, the gibberish and some arm flapping were all part of a serious lesson to help kids learn the telltale signs of a stroke by imitating them. The idea is to enlist children, particularly those who may live with older relatives, as an army of eyes to help recognize the warning signs, get help for victims more quickly and hopefully save lives. “If my mom has a stroke, I’ll know what to do,” said 10-year-old Madison Montes. “Run to the phone and call 911.” The experimental health education program at Montefiore Medical Center in the Bronx is aimed at the most crucial factor when it comes to a stroke: time. Each year, about 795,000 Americans have a stroke and about 130,000 die. Some are caused by bleeding in the brain, but the vast majority is caused by a clot that blocks blood flow, starving brain cells. The drug TPA can dissolve those clots and reduce disability but only if it’s given within three to four hours of the first symptoms, and the sooner the better. Yet only about 5 percent of patients receive it, in part because many stroke sufferers don’t get to the hospital in time for testing to tell if they’re a candidate. The early warning signs include a droopy side of the face, slurred or strange speech, and the inability to keep arms raised. “There’s a pretty good chance some children might witness a parent or a grandparent having a stroke,” said Jim Baranski, CEO of the National Stroke Association. “So if they’re armed with the signs and symptoms, they could likely save a life.” Montefiore’s program, one of a handful tried across the country, has been used since 2012 with private schools in its neighborhood, where children are often in a grandparent’s care because parents are absent or both working. The goal is to study the results and, if successful, replicate the program across the country. “The kids get a kick out of it because they get to do a little acting,” said Dr. Robert Glover, a neurologist who helped develop the program. “But when they’re done, they know about stroke and they can teach their families.” Dr. Kathryn Kirchoff-Torres, who led the class from St. Ann’s School in the Bronx, said the kids are already “little message machines” bringing home from school what they learn about the benefits of exercise, not smoking and eating well. At the start of the stroke class last month, in a first-floor room at the hospital, the doctor asked, “Who knows what a stroke is?” “A heart attack?” one child offered. “Well, we like to call it a brain attack,” Kirchoff-Torres said. “It’s a problem with the brain.” She then taught the children to use the word “FAST” as a memory device. With cartoons and music bringing the point home, they learned “F” is for face, “A” is for arms, “S” is for speech and “T” is for time. After the play-acting and the multimedia show, the doctor invited questions from the children. “How do you catch a stroke?” one boy asked. The doctor assured him that strokes are not contagious but can be caused by “high blood pressure, smoking cigarettes, junk food.” “What if we don’t have a phone?” a girl asked. Kirchoff suggested asking a neighbor or running to a storefront. “What if you live in the desert?” was the follow-up question, to which Kirchoff smiled and said, “It’s a good thing you live in the Bronx.” After the class, the children were presented with pens labeled “FAST” and with pencil erasers in the shape of human brains, which were very popular. One parent in attendance, Jason Sawtelle, said he felt the lesson “plants the seed” “Maybe not every 8-year-old is capable of this,” he said, “but some certainly are.” Similar programs have been tried elsewhere. At Harlem Hospital in Manhattan, Dr. Olajide Williams uses hip-hop music to engage New York City fourth- through sixth-graders in three hours a year of stroke education. “Beyond sixth grade, the kids become a little too cool,” Williams said. But with the younger kids, “we’ve shown that these children can learn basic stroke pathology.” He said at least two children have been credited with helping to save lives. A similar stroke education study in the Corpus Christi, Texas, public schools used Tejano music. It was headed by Dr. Lewis Morgenstern, director of the stroke program at the University of Michigan Medical School. The 2007 study found that the youngsters responded well. “The data was highly positive in terms of knowledge about stroke and their intention to call 911,” Morgenstern said. “The earlier we can make people aware of stroke and that it’s arguably the most treatable of all catastrophic conditions, the better off we will be.”

Augmented reality combines real world attributes, such as physical position or movement, with computer-generated, often graphical, data. Together, the attributes and data form a user experience occurring in a real world environment, while interacting with virtual objects, typically occurring in real time. In general, augmented realities are presented through graphical user interfaces, which can also be supplemented with auditory cues, tactile feedback, and other computer outputs. Thus, the user is able to view virtual objects positioned in the real world. Recently, the capabilities of mobile platforms, such as portable media players, personal data assistants (PDA), and mobile telephones, have grown to enable on-the-go generation of augmented realities. Users are no longer tied to desktop environments and the physical constraints imposed by having to remain in a relatively stationary position. Moreover, these platforms are increasingly being equipped with input capabilities extending beyond the immediate user interface. For instance, mobile telephones often incorporate built-in cameras and global positioning system (GPS) receivers, the latter of which enables the device to determine the user's physical location. Similarly, portable media players are beginning to incorporate motion sensors that sense directional device orientation. For the most part, the extended capabilities of these devices have been limited to providing abstract objects with a real world view for a single user, rather than collaboratively between multiple users. For instance, Pokemon Go by Niantic, Inc. of San Francisco, Calif., is a location based augmented reality game in which a user captures virtual Pokemon characters by moving around in the real world. A user can capture the characters by touching a virtual ball that appears on a screen of a mobile device and flinging the ball towards the virtual character positioned in the real world. However, catching the virtual characters are limited to a single user who is on a solo journey to catch and collect Pokemon characters. In a recent update, raid battles were introduced in Pokemon Go, in which multiple players can join a raid to collaboratively attack a raid boss by independently fighting the raid boss at the same time. However, within the raid, a user is unable to interact with other users in the same battle or raid. For instance, the users are unable to communicate with one another and to coordinate moves to defeat the raid boss. Further, views of the users are not coordinated such that cooperative play can occur. Therefore, a need remains for an approach to providing spatially accurate and oriented perspective views of virtual objects in the real world based on the location of each of a plurality of users.

Pumpkin Seeds: Benefits, Nutrition and Diet Tips Benefits of Pumpkin Seeds Pumpkin seeds also known as pepitas are edible seeds of the pumpkin. These seeds are flat and symmetrically oval. They are encased in a white-colored outer husk and their inner seed is green in color. They are subtly sweet and have a nutty flavor. They are native to America Pumpkin seeds are loaded with numerous micro nutrients, they are rich in a healthy fat like oleic acid and linoleic acid. Pumpkin seeds are widely used in Mexican cuisine, they are usually roasted and used as a snack. They are packed with energy and nutrition. They are the only seeds which are alkaline nature. Following are the few Health Benefits of Pumpkin Seeds • Anthelmintic Effect Pumpkin seeds are widely used as a natural remedy for parasites in many cultures. Few Native American tribes use pumpkin seeds to get rid of tapeworms in the stomach. • Rich in Magnesium Magnesium – have a role in 300 – 500 functions in the human body, it is needed in trace amount but participate in many functions. • Rich in Antioxidants Pumpkin seeds are rich in diverse antioxidants, like- contain vitamin E, alpha-tocopherol, delta-tocopherol, alpha to comonoenol, and gamma-tocomonoenol, hydroxybenzoic, etc. hence pumpkin seeds helps in protecting our body from free radicals which are naturally generated in the body by stress, sun exposure and as part of the normal aging process. Pumpkin seeds also help to heal of cancer like prostate cancer and bladder and colon and breast cancer. • Anti-inflammatory and antimicrobial Pumpkin seeds are one of the most recommended natural anti-inflammatory foods. Pumpkin seed oil is also recommended for its natural antimicrobial properties. The unique proteins in pumpkin play a major contribution in antimicrobial action of pumpkin seeds. Pumpkin seeds are a natural source of protein which is almost 25 – 50 % of the total weight of the seeds. • Maintains Blood pressure Pumpkin seeds help in lowering blood pressure, they play a vital role in relaxing the blood vessels. High magnesium in pumpkin seeds is effective in reducing the risk of a heart attack. Omega 3 and omega 6 fatty acids present in pumpkin seeds help in lowering cholesterol in the blood. They also decrease the risk of thrombosis and arrhythmia’s. Likewise, pumpkin seeds are also beneficial in reducing atherosclerosis. Maintains Blood pressure • Heart health Studies have depicted that pumpkin seeds may help in reducing cholesterol. A substantial amount of research has been conducted which endorse the hypolipidaemic effect of pumpkin seeds. Pumpkin seeds help in a significant decrease in lipid parameters. • Healthy Bones Pumpkin seeds a good source of minerals like calcium and magnesium, which plays a vital role in decreasing osteoporosis. Pumpkin seeds contain a considerable amount of zinc, which is important for bone health. • Maintains blood Sugar Level While Diabetes patients are advised not to consume pumpkin because of high carbohydrate content, the seeds are not barred from consumption. • Male Infertility Pumpkin seeds are rich in zinc, which is an essential mineral it helps in enhancing the Sperm quality of man. It helps in increasing the sperm counts and enhances the motility of sperm. Zinc is also vital for prostate gland function. It also helps in the growth of reproductive organs during development. Likewise, zinc is a vital fertility nutrient, it helps in improving the sexual health status of both male and female. • Enhances sleep quality Sleep is an important ingredient in health. Zinc is required for the production of tryptophan, it is an important amino acid which promotes sleep. It is converted into serotonin which converts into melatonin. Melatonin plays an important role in the regulation of sleep and human circadian rhythms. It is a sleep facilitator in human thus helps in the treatment of insomnia. Enhances sleep quality Pumpkin seeds are also beneficial in decreasing the symptoms associated with Leg cramps, calf pains, and migraines. It is always better to eat soaked pumpkin seeds, about 1 to 2 tablespoon is sufficient for a day. Benefits of Pumpkin Seeds Pumpkin seeds also known as pepitas are edible seeds of the pumpkin. These seeds are flat and symmetrically oval. They are encased in a white-colored outer husk and their inner seed is green in color. They are subtly sweet and have a nutty flavor. They are&hellip; Pumpkin Seeds: Benefits, Nutrition and Diet Tips Pumpkin Seeds: Benefits, Nutrition and Diet Tips 2019-08-21 Dr. Reena Prajapat Rate this article (click stars to rate) Audience / Patient: Health Professionals: Beauty Professionals: 100 Amazing Dr Reena Prajapat is a dentist, she post graduated in hospital management from Apollo Institute of hospital Administration, Hyderabad. Her strong belief in natural and home remedies lead her to Pursue bachelor in naturopathy.

Louisiana State University entomologist Rodrigo Diaz said researchers only recently discovered the foreign family of insects to which the invasive species belongs, called Aclerdidae, which is native to Japan and China. But the lab tests that identified it couldn’t reveal how it arrived — on a ship, attached to a migrating bird or even on the wind. What’s certain is that last year a team of surveyors checking the cane, which comes in various shades of green, found stalks bent in water, brown and dead in the mouth of the Mississippi River. That’s when they began to notice significant die-offs of four varieties of cane and more open water in the hundreds of acres the cane once occupied. Two to three years before, there was a thick, unrelenting wall of marsh. AD AD “They are feeding on it,” said J. Andrew Nyman, a professor at the School of Renewable Natural Resources at LSU. “The bugs suck the sap out. The leaves are trying to send sugar to the roots, and they suck out so much that the plant can’t function. It dies.” Nyman said he picked up a stalk one day, and the bugs covered it. In an online seminar about the invasion, Diaz displayed a graph that said six feet of cane can have nearly 200 insects almost invisible to the naked eye, and 700 in extreme cases. He described them as translucent, extremely small but highly mobile, and moving in dense populations. The mealy bug, as it’s also known, adds to the many invasive insects that have made their way to the United States from Asia and Europe. Animals that have ruined plants and crops, and even invaded homes, include the brown marmorated stink bug, the kudzu bug, the emerald ash borer, and aphids that are threatening Florida’s citrus industry by destroying orange groves and wreaking havoc on Christmas tree pines from West Virginia to North Carolina along the Appalachian Mountains. AD AD The Louisiana cane is crucial to staving off land loss. It builds soil in an area that lost 250 square miles of coast to erosion and sinking land over about a half century. Roseau cane also has a checkered history. At least one variety, the European, is itself invasive. The delta Roseau cane from North Africa and the greeny variety from Europe were introduced. Scientists aren’t sure whether a fourth variety, the Gulf Roseau cane, is invasive or introduced. But on a coast beaten down by harsh weather and erosion, even the exotic species turned out to be a godsend for stabilizing the land. In other areas of Louisiana and other states, they are a nuisance, which is why some people, Diaz said, want to transport the parasite to those areas to destroy the cane. Bad idea, he said. Who knows where the bugs will go after the cane is gone. They could wander to native plants, even farm crops, and develop a taste for those. But the insects don’t wait for human transport. Even now they hop aboard grackles and red-winged blackbirds that flock into the reedy cane, hoping for a ride to another healthy batch of cane to get a meal and lay eggs. They’re also known to hop into air boats piloted by fishermen. Scientists have pleaded with the owners to wash them down. Diaz said experts are weighing options to fight the threat. An idea for a controlled burn is derived from China, where blazes are set in marshes to get rid of the swarm. But fire can spread, and the Louisiana coast has a network of oil and gas wells that could explode in flames. AD AD Planting a plague-resistant strain of cane has come up, but no one knows whether such a thing exists. Insecticide is another consideration, but that could get expensive because of the vast area that needs spraying and would require environmental assessments. Another concern is the inadvertent contamination of speckled trout, redfish, shrimp and oysters — not to mention the birds that feed on them. Maybe the parasite’s natural predator, a tiny wasp from Asia, can be employed. They have been considered as a solution to the stink bug problem as well, but it takes years of study before a new species can be released, for fear that the wasp might turn on native animals and start killing them, too. Or maybe, just maybe, the insect is its own worst enemy. It might eat itself out of house and home before anything else can be done. “What you could end up with in 10 years is the bugs will die back,” Nyman said, after killing all the cane. “They’ll become a smaller part of the landscape, both the bug and the plant.”

Synovitis is the medical term for inflammation of the synovial membrane. This membrane lines joints which possess cavities, known as synovial joints. The condition is usually painful, particularly when the joint is moved. The joint usually swells due to synovial fluid collection. In principal, changes of the synovial membrane can be inflammatory or non-inflammatory. To the latter group belong some benign tumors, such as tenosynovial giant cell tumor, lipoma or synovial chondromatosis. Rare non-inflammatory changes are the group of storage diseases. Inflammatory synovial diseases can be differentiated into crystal-induced arthropathy, such as gout and pseudogout, granulomatous diseases, such as tuberculosis, sarcoidosis and foreign body reactions and into the large group of non-granulomatous synovitis. This last group is by far the most common and often causes difficulties in assigning the histopathological findings to a definite diagnosis. Hence, synovitis may occur in association with arthritis as well as lupus, gout, and other conditions. Synovitis is more pronounced in rheumatoid arthritis than in other forms of arthritis, and can thus serve as a distinguishing factor, although it is also present in many joints affected with osteoarthritis (OA). Visco-supplementation is the process of injecting a gel-like substance into the joint. The substance is thought of as an additive to the joint fluid, thus lubricating the cartilage, and improving joint flexibility. This method of treatment, however, requires ongoing injections, as benefits are only temporary, because the currently used substances are degradable within weeks to months. Substances used in visco-supplementation include hyaluronic acid, or HA (Legend®, Hylartin® and Synacid®, Synvisc, Euflexxa, Supartz etc) and poly-sulfated glycosaminoglycans (PSGAGS) such as Adequan®. Polyacrylamide hydrogel (PAAG) is known for its ability to support cellular growth in vitro and allow in-growth of host tissue cells in vivo. This has been documented in the soft subcutaneous tissues of mice, rats, rabbits, pigs and humans. The tissue integration begins immediately after PAAG injection as a sort of foreign body reaction to the PAAG. Host macrophages and foreign-body giant cells initially surround the PAAG and then invade it. In the process these cells are gradually transformed into fibroblasts and endothelial cells, which eventually form a thin vessel-bearing fibrous network inside the PAAG. However, the integration of PAAG in these tissues was not associated with any luminal surfaces and the PAAG contained macrophages and giant cells up until 14 months post injection. WO 02/16453 discloses the use of a polyacrylamide hydrogel (PAAG) for treating e.g. arthritis. where the treatment is considered to be based on a lubricating and cushioning effect of the hydrogel. WO 2012/123385 discloses use of PAAG in the treatment and/or prevention of joint swelling and/or bone oedema in a mammal suffering from arthritis. Neither WO 02/16453 nor WO 2012/123385 discloses prevention and/or treatment of synovitis in a mammal. Christensen and Daugaard (J Arthritis 5: 217; Sep. 16, 2016) provide a case report on the histological appearance of the synovial membrane after treatment of knee osteoarthritis with polyacrylamide gel injections. The common symptoms for all types of synotivis in mammals include varied levels of pain in mammals. Synovitis symptoms can be treated with anti-inflammatory drugs such as NSAIDs. Another possibility is injection of steroids directly into the affected joint. Specific treatment depends on the underlying cause of the synovitis. Common to the various types of existing treatment is that they all have their disadvantages, e.g. short term treatment, toxicity and side-effects. Hence, there is a need for an alternative or improved way to prevent and/or treat synovitis in a mammal. Also, there is a need for an alternative or improved way to prevent and/or treat synovitis pain in a mammal. Surprisingly, the present inventors have found that polyacrylamide hydrogel (PAAG) is useful in the prevention and/or treatment of synovitis. Without being bound by a particular theory, the analgesic effect on synovitis is considered caused by a stable, long-lasting sub-synovial layer of PAAG traversed with thin strands of connective tissue and changes to synovial cell composition or cytokine production. Accordingly, the present invention surprisingly demonstrated the formation of a novel synovial lining layer after integration of the PAAG into mammal joints that persisted for at least 24 months.

NIH researchers have recently announced the isolation of a human B-lymphotropic virus (HBLV) from six individuals with various lymphoproliferative and immune disorders. HBLV appears to belong to the group of human herpes viruses and may be associated with B-lymphoproliferative disorders. Serological tests using indirect immunofluorescence methods revealed that only 4 of 220 random healthy blood donors were seropositive for HBLV antibodies. An ELISA-based assay for HBLV antibody detection would facilitate further study of the prevalence of infection with this virus and the relation of infection to human disease. So far, however, HBLV has not been seen to be infectious for any cells other than fresh B-lymphocytes. In this proposal, the investigators present their plan for using HBLV cultured in B-lymphocytes obtained from umbilical cords to demonstrate the feasibility of developing an ELISA-based test kit for HBLV antibody detection. They also discuss their plans for establishing a cell culture system other than fresh B-lymphocytes for virus proliferation, so that tests can be produced on a commercial scale.

We use cookies to enhance your experience on our website. By continuing to use our website, you are agreeing to our use of cookies. You can change your cookie settings at any time.ContinueFind out more In numbers after 30 the conjunction y is used between the tens and the units, but not between the hundreds and the tens: 46 = cuarenta y seis 432 = cuatrocientos treinta y dos millón requires the use of de: a million people = un millón de personas three million votes = tres millones de votos Numbers over one million 1 000 000 000 = mil millones or un millardo (one billion) 1 000 000 000 000 = un billón, un millón de millones (one trillion) 1015 = mil billones (one quadrillion) 1018 = un trillón (one quintillion) The examples above show Spanish usage as recommended by the RAE. However, billón and trillón are increasingly being used as the equivalents of the English billion and trillion respectively. This usage is not sanctioned by the RAE. It should be noted that the Real Academia Española (RAE) recommends the International Standardization Office's ruling on the treatment of thousands. This states that spaces should be used to separate groups of three digits. This has been followed in the list above. However, in practice it will be found that spaces, points, and commas are used as separators. In most Spanish-speaking countries a point is used for writing figures over one thousand: 1,000 (one thousand)= 1.000 (mil) 1,000,000 (one million)= 1.000.000 (un millón) Some Latin American countries, however, use the comma as in English: 1,000 (one thousand)= 1,000 (mil) 1,000,000 (one million)= 1,000,000 (un millón) Gender and agreement Numbers in Spanish are masculine when used as nouns, and require an article:

Recent Articles At times it might seem to some of us as though the world's top boffins are slacking at their task of making our technology better and more advanced: but not today. Today we learn that some of them are on the track of something which everyone involved in IT must have been lusting after for years. Wouldn't it be nice, you must have thought, if instead of all this plastic and silicon and glass and so forth, my kit - displays, chips, tablets, laptops, all of it - could be made instead of lovely concrete? Well, maybe it can. Topflight brainboxes at the US government's Argonne National Lab have been toiling away on this for some time, and they believe they may have cracked a vital enabling tech trick - to wit, that of getting liquid cement to behave like a metal. According to an Argonne statement: This makes cement a semi-conductor and opens up its use in the profitable consumer electronics marketplace for thin films, protective coatings, and computer chips. Not only the chips and the casing of your device could be made of the rough grey material of tomorrow, but even the display. “This new material has lots of applications, including as thin-film resistors used in liquid-crystal displays, basically the flat panel computer monitor that you are probably reading this from at the moment,” explains Chris Benmore, Argonne Lab physicist. In essence the crafty trick perfected by Benmore and colleagues around the world is that of getting cement to become a "metallic glass", aka "liquid metal". Such materials have long been experimented on in labs, but so far it's only been possible to make them out of metal and they haven't yet taken the world by storm as many have thought they might - though the SIM removal tool in some early iPhones and 'Pads was made of such stuff. “This phenomenon of trapping electrons and turning liquid cement into liquid metal was found recently, but not explained in detail until now,” Benmore says. “Now that we know the conditions needed to create trapped electrons in materials we can develop and test other materials to find out if we can make them conduct electricity in this way.” The problem with using metallic glasses in the real world is that they need to made very hot to form, so hot that moulding and handling them becomes very difficult. Benmore and his colleagues overcame this with their test batches of liquid metal cement by using an "aerodynamic levitator with carbon dioxide laser beam heating" - though they don't specify (or know yet) whether this method would be suitable for industrial use. All in all, the new super metallo cement glass stuff would seem to be a lot further from actual use than regular metalglass/liquidmetal - and even that seems to be somewhat bogged down. But the prospect of concrete* computing seems intriguing nonetheless. ® Bootnote *We do know that cement isn't the same as concrete. For those that don't, cement is mixed with "aggregate" - other stuff such as broken rock, pebbles or whatever else may be handy - to make concrete. Cement is basically the glue that holds concrete together. What the boffins are on about above is cement. However should one in future make a display, chip or whatever using liquid-metallic cement tech, the finished product would surely also include other things just as concrete includes other things than cement. So it might seem fair enough to call it a "concrete computer".

Q: Why shift the array in the &max subroutine example in Learning Perl, 6th ed., ch. 4 Why is the array shifted at the beginning of the subroutine? sub max { my($max_so_far) = shift @_; foreach (@_) { if ($_ > $max_so_far) { $max_so_far = $_; } } $max_so_far; } Is it just to give $max_so_far an initial value? The program runs exactly the same with my($max_so_far) = undef; Is there a particular reason to shift the array to start with? (I ask because I spent about 10 min trying to figure out why that shift was essential to the subroutine.) A: The program does not run exactly the same if you initialize $max_so_far to undef. What if all of the input values are negative numbers? A: It's to initialise the return value for the function and something the rest of the function can compare against. Consider some scenarios and trace through the code. For example, say it is invoked like this: my $max = max(1,2,3); Inside max, the first line sets $max_so_far to 1 and @_ becomes (2,3). Now when we run through the foreach loop we have an initial value and avoid undef errors. It first compares $max_so_far to 2, updates it to 2, and so on. Another example is if max is invoked like this: my $max = max(1); Inside max, the first line sets $max_so_far to 1 and @_ becomes (). When we hit the foreach loop, it has nothing to iterate through and just returns the initial value $max_so_far. Another example is if max is invoked like this: my $max = max(); Inside max, the first line sets $max_so_far to undef because @_ is empty. There's nothing to iterate through in the foreach loop, so the function just returns undef.

Q: Can you really redefine keywords in C language? Can you explain the code below? How we can use #define for a keyword of C? #include <stdio.h> #define int int* int main(void) { int *p; int q; p = 10; q = 5; printf("%d %d", p, q); // your code goes here return 0; } Output: 10 5 A: This #define int int* is a preprocessor macro. If you want to define your own synonyms for types then use typedef. You cannot create keywords in a language you did not create. Sample: #include <stdio.h> typedef int * myIntPtr; int main(void) { int i = 10; myIntPtr x = &i; printf("%d", *x); return 0; } Output: 10 Also semantically making an int to an int * makes no sense.

Background ========== The use of ecological models to describe and understand physical activity and other health behaviors calls for new strategies for assessing and intervening upon the social and physical environments that affect those health behaviors \[[@B1]\]. While most attention has been paid to assessing activity options in the neighborhood environment \[[@B2]-[@B5]\], less attention has been focused on the home environment. The home environment, including areas both inside and outside (e.g., yard, garden) as well as items in the home that either facilitate or discourage activity, may have an important influence on physical activity patterning and warrants further study. Valid and reliable methods of assessing environmental supports for physical activity and competing sedentary media pursuits within the home environment are needed. Assessments of both the activity and sedentary environmental supports are warranted since physical activity and sedentary behavior are unique and independent behaviors \[[@B6],[@B7]\]. A person can be physically active but also engage in substantial amounts of sedentary activities like watching television or playing video games \[[@B7]\]. Particularly in children, there is good evidence that television time is related to the prevalence of childhood obesity \[[@B8],[@B9]\]. To date, most physical activity inventories have shown weak (r = 0.04 to 0.19) \[[@B10]\] or no association \[[@B11]\] between physical activity and children\'s perception of their access to sporting and fitness equipment at home. In a study of college students, the quantity of home exercise equipment based on dichotomous (yes/no) response checklists was only associated with strength training activity in bivariate analyses (r = 0.18) \[[@B12]\]. In another study of college students, a similar checklist was also associated with strength exercise, not total physical activity, using multivariate analyses (p = 0.005) \[[@B13]\]. These low associations may suggest that equipment is not related to activity or may reflect methodological issues. A more comprehensive assessment of the home environment may help to better explain physical activity and sedentary behaviors and the specific activities being performed. Other recent work by Hume et al. \[[@B14]\] has sought to develop a similar checklist-based tool to assess the home environment for physical activity among Australian children. Evaluation suggested that test-retest reliability of this instrument was high for most items (Kappa = 0.85 to 1.00). However, the validity of the tool was not examined, and thus this tool may suffer similar shortcomings as those checklists previously described. While the checklist-based instruments assess the breadth of items (the number of unique items), a more comprehensive assessment would also measure magnitude (how many of each item are present), a characteristic of the home environment which may provide an important influence on behavior. In addition, previous home environment assessment tools have focused on availability, not accessibility. Accessibility is related to ease of use and cueing of behavior. A toy or piece of sports equipment that is readily available and easy to access poses a smaller barrier to use than one that must be retrieved from some inconvenient space. Likewise, the presence of a television set in the bedroom or a basketball hoop in the yard may serve as a visual cue to make use of those pieces of equipment. Similarly, stimulus control (i.e. controlling stimuli that can encourage or discourage a behavior) is recognized as one of the important processes of change used in conjunction with the Transtheoretical Model of behavior change \[[@B15]\] and has been associated with change from contemplation to preparation stages \[[@B16]\]. An inventory that can help evaluate access may be an important contribution to measurement tools. The purpose of this study was to develop and test the reliability and validity of a self-report survey instrument to comprehensively reflect the availability and accessibility of physical activity and media equipment in the home environment. It is our hope that such an instrument may help us better understand how the environment impacts the health and health behaviors of families. Methods ======= This research was conducted as part of the Identifying Determinants of Eating and Activity (IDEA) study for the purpose of designing and assessing the reliability and validity of an instrument to assess the home environment with regard to equipment related to physical activity and screen time. IDEA is an etiological study to evaluate potential determinants of the development of overweight and obesity in youth; it is being conducted in the metropolitan area of Minneapolis and St. Paul. IDEA assesses potential obesogenic factors in youth\'s personal, family, home, school, and neighborhood environments using a longitudinal cohort design (Lytle, in review). Physical Activity and Media Inventory (PAMI): instrument development -------------------------------------------------------------------- The PAMI was designed to be a self-report inventory of both the availability and accessibility of equipment and other resources that may support family members\' participation in activity and sedentary behaviors. We were interested in creating an instrument that documented the physical environment, not individual behavior within the environment, therefore, we chose not to assess the frequency with which equipment inventoried was used. The development of the PAMI began by examining and adapting existing inventory instruments to address the specific needs of our research objectives and by creating a draft inventory. This measurement tool went through multiple drafts where face validity, clarity of the questions and format, and feasibility of administration was reviewed by study investigators and with other researchers who study family health. When a satisfactory draft was complete, we asked a convenience sample of ten families, recruited through local park and recreation centers, to complete and comment on the inventory. All participants were able to complete the PAMI and feedback suggested that the PAMI would be feasible as a self-report instrument. Based on their feedback, we further refined the instrument to improve instruction clarity and form layout; this version was used to assess the reliability and validity of the PAMI instrument. The PAMI version that was evaluated in this reliability/validity study included a list of 50 physical activity equipment items grouped by the following categories: sports equipment, fitness equipment, transportation equipment, foot wear, water sports, and outdoor/yard (see Additional File [1](#S1){ref-type="supplementary-material"}). There also were five media equipment items listed: television, video cassette recorder (VCR) and/or DVD, digital video recorder (DVR) and/or TiVo, video game system, computer (desktop or laptop). To obtain more information on media, six media-related questions were asked including: 1\. Type of television programming available (No TV, No cable, Basic Cable, Cable + Premium channels, Satellite/Dish), 2\. Number of channels received (No TV, \<15, 15--30, 31--45, 46--60, \>60), 3\. Number of videos and DVDs in the home (0, 1--25, 26--50, 51--75, 76--100, \>100), 4\. Number of video and computer games in the home (0, 1--10, 11--20, 21--30, 31--40, 41--50, \>50), 5\. Type of internet service (No internet service, Dial-up modem, DSL Modem, Cable Modem, Don\'t know), 6\. Size of the primary television in the home (diagonal screen size, in inches). Opposite the list of items there was a list of 16 rooms (e.g., living room, kitchen, etc\...) which also included Yard/Outdoor Space, Attic/Basement/Storage Area, Garage, and Automobile(s) where respondents indicated the location of the available activity or media equipment. There was space for three \"other\" rooms in case the home had a room not listed (e.g., sunroom) or more space was needed for a particular location (most common for the garage). To determine validity, one parent/guardian from each participating family completed the PAMI while a trained research assistant simultaneously completed an inventory in the family\'s home. To determine reliability, the same parent/guardian completed a second inventory one week after the first administration. The second inventory was done without the presence of a research assistant. All data were collected during May-June, 2006. Subjects -------- Families were recruited to participate in this study through posted flyers and staffed information booths at four Minneapolis Park and Recreation Department buildings. Eligibility criteria included that each family had at least one child between the ages of 10 and 18 and an adult (parent/guardian) willing to complete the inventory and permit a home visit for validation purposes. Written informed consent was obtained prior to all data collection activities. The University of Minnesota Institutional Review Board approved this study. Procedures ---------- At the first visit, one adult household member and a research staff member did a walk through of the house, independently completing the PAMI instrument. The respondents were asked to walk through each room and indicate which physical activity and media equipment items were present (using code numbers from the list) and to rate the accessibility of each item. Access response options were: put away and difficult to get to; put away and easy to get to; in plain view and difficult to get to; in plain view and easy to get to. The research assistant independently completed the criterion PAMI at the same time with as little communication with the participating family member as possible. If a participant began speaking, the research assistant would ask them not to speak. Due to the potentially intrusive nature of such an inventory, the research assistant only inventoried those items that were in plain view, without moving or looking under furniture or opening closets, unless the participant did so first. At the end of the first visit, the research assistant left a blank PAMI instrument with the participant. The participant was asked to independently complete their second PAMI approximately one week later using the same approach as the first administration. A self-addressed stamped envelope was provided for mailing the completed form back to research staff. Reminder messages (telephone and email) were made to families one day prior to the day when the second inventory was to occur. Demographics were collected by self report during the first PAMI administration and included the age, race, and education level for each member of the household and one question asking about their homeowner status (i.e., apartment, condominium, multi-family house, or single-family house). Data reduction -------------- The PAMI data were reduced to the following primary variables, calculated separately for physical activity and media equipment: total number of items, the density of items in the home (total number of items divided by the total number of rooms/locations), the number of items located in bedrooms, the density of items in bedrooms, the number of televisions in the home, and the number of televisions in bedrooms. For comparability to a checklist type of instrument \[[@B12],[@B13]\], which record only the presence or absence or an item and not the quantity, the number of unique items reported was separately summed for physical activity equipment. For example, if more than one bicycle was reported on the PAMI, bicycle was only counted once to most closely represent a checklist type instrument. We also created and assessed two summary scores. First, we calculated separate summary scores that accounted for availability and accessibility of the physical activity equipment (Physical activity Availability and Accessibility Summary Score (PAASS)) and media equipment (Media Availability and Accessibility Summary Score (MAASS)). Each item was multiplied by the accessibility code with 1 = \"put away and difficult to get to\" and 4 = \"in plain view and easy to get to\". A higher score reflects a greater overall presence in the home (both availability and accessibility). To provide more detail, we examined specific categories of items, including; sports equipment, fitness equipment, transportation equipment, water sports equipment, and outdoor/yard equipment). Within each of these categories, we calculated the total number of items and mean accessibility ratings for all items within each category. To rank the overall quality of the home, an overall home environment score was also calculated as the ratio of the PAASS to the MAASS (referred to as the Activity:Media Ratio Score). A higher overall Activity:Media Ratio Score would reflect a home more conducive for being physically active and less sedentary. Statistical analysis -------------------- SAS version 9.1 (Cary, NC) was used for all data analyses. A p ≤ 0.05 significance level was used as a guide for identifying significant relationships. Test-retest reliability of the continuous variables from the PAMI (i.e., number and density of items, number and density of items in bedrooms, checklist quantity, summary scores, number of physical activity items within categories, and overall home environment score) was assessed by intraclass correlation coefficient (ICC, 95% confidence intervals). Test-retest of the accessibility ratings (overall for the home and by physical activity item categories) was assessed by Mantel-Hanzel chi square analyses. The reliability of the additional media-related questions was assessed by percent agreement (Kappa coefficients, 95% Confidence Intervals). Validity was evaluated by examining agreement between data from the participants and research assistants using Pearson correlation coefficients for continuous variables. Mean differences in continuous variables between the participant and the research assistant were identified with two-tailed independent t-tests. Comparison of the accessibility ratings between the participant and research assistant was calculated by Mantel-Hanzel chi square analyses. Results ======= A total of 31 families agreed to participate in the study and data from all families are included in these analyses. All participating families completed the PAMI at least once for validation purposes, and over 77% also completed the PAMI a week later to assess reliability. The PAMI took, on average, 40 minutes to complete but was dependent on the size of the home and the number of items present. Table [1](#T1){ref-type="table"} provides the demographic characteristics of the participating families for the validation (N = 31) and reliability (N = 24) samples. Dependent t-tests and chi square analyses revealed no demographic differences in the families included in the validity and reliability samples (p = 0.26--0.94). The majority of parents/guardians were female, Caucasian and had at least a college education. On average, there were approximately four people per home, most with two adults and two children living in single-family homes. ###### Subject characteristics for the validity and reliability samples **Validity Sample** **Reliability Sample** ------------------------------------------------- --------------------- ------------------------ **Age, mean ± SD** 38 ± 7.2 37 ± 6.6 **Race/Ethnicity, n (%)**  Caucasian 16 (52%) 12 (50%)  African American 6 (19%) 6 (25%)  Mexican-American, Puerto Rican, Latin American 2 (6%) 1 (4.2%)  Native American 2 (6%) 1 (4.2%)  Asian American 2 (6%) 1 (4.2%)  Other/Unknown 3 (10%) 3 (13%) **Highest Education, n (%)**  Did not graduate high school 2 (6%) 1 (4%)  High school/Trade school/Some college 7 (22%) 5 (21%)  \> College graduate 20 (65%) 14 (58%)  Not answered 2 (6%) 4 (17%) **Number of people in home, mean (SD)** 3.9 (0.9) 3.8 (1.1%) **Number of adults in home (≥ 18 yrs), n (%)**  1 Adult 10 (32%) 10 (42%)  2 Adults 20 (65%) 13 (54%)  3 or More Adults 1 (3%) 1 (4%) **Number of minors in home (\< 18 yrs), n (%)**  1 Child 5 (16%) 4 (17)  2 Children 17 (55%) 13 (54)  3 Children 6 (19%) 5 (21)  4 Children 3 (10%) 2 (8) **Type of home, n (%)**  Apartment 2 (6%) 2 (8%)  Multi-family house 5 (16%) 5 (21%)  Single-family house 23 (74%) 15 (63%)  Not answered 1 (3%) 2 (8%) Reliability ----------- Table [2](#T2){ref-type="table"} contains the results of the reliability analyses. Test-retest for the physical activity and media equipment variables was good to excellent (ICC = 0.71 to 0.96). Use of a checklist approach results in high reliability for the physical activity equipment (ICC = 0.93). The proportion of physical activity and media equipment items in each accessibility category was similar between the first and second administration (Mantel-Hanzel chi square tests, p = 0.50 and 0.30, respectively). The calculated PAASS and MAASS (number of items \* accessibility) revealed high reliability between the first and second administration of the PAMI (ICC= 0.87 for physical activity items and ICC = 0.93 for media equipment). ###### One-week test-retest reliability of PAMI variables; mean ± SD with Intraclass Correlation Coefficient (ICC) and 95% Confidence Interval (CI) or percent with Mantel-Hanzel Chi Square **Variable** **Time 1** **Time 2** **ICC (95% CI)** ----------------------------------- ---------------- ---------------- -------------------------- Total number of rooms 11.9 ± 3.10 12.5 ± 3.22 0.86 (0.70 -- 0.94) Physical Activity Equipment  \# of Items 65.7 ± 56.52 64.8 ± 58.06 0.94 (0.87 -- 0.98)  Household Density 5.1 ± 4.16 5.0 ± 4.53 0.84 (0.86 -- 0.97)  Ave \# of Items/Bedroom 3.3 ± 3.10 3.5 ± 3.30 0.76 (0.51 -- 0.89)  Checklist Quantity 16.4 ± 7.85 17.5 ± 7.63 0.93 (0.84 -- 0.97)  Accessibility   Away and difficult to get to 8.7 4.4 M-H χ^2^= 0.46, p = 0.50   Away and easy to get to 17.8 21.5   In view and difficult to get to 2.9 2.3   In view and easy to get to 70.3 71.9  PAASS 207.4 ± 166.74 216.8 ± 190.12 0.87 (0.74 -- 0.94) Media Equipment  \# of Items 8.4 ± 5.02 8.4 ± 4.93 0.96 (0.91 -- 0.98)  Household Density 0.73 ± 0.38 0.73 ± 0.43 0.89 (0.76 -- 0.95)  Ave \# of Items/Bedroom 2.4 ± 0.78 2.1 ± 0.68 0.72 (0.44 -- 0.87)  \# of Televisions 2.5 ± 1.20 2.7 ± 1.27 0.89 (0.76 -- 0.95)  \# of Televisions in Bedrooms 1.17 ± 0.94 1.27 ± 1.03 0.93 (0.82 -- 0.97)  Accessibility   Away and difficult to get to 1.4 2.6 M-H χ^2^= 1.05, p = 0.30   Away and easy to get to 10.3 12.7   In view and difficult to get to 0.7 1.3   In view and easy to get to 87.6 83.4  MAASS 32.0 ± 25.5 31.2 ± 18.84 0.93 (0.84 -- 0.97) PAASS, Physical activity Availability and Accessibility Summary Score MAASS, Media Availability and Accessibility Summary Score The test-retest reliability for the number of items within all physical activity equipment categories was high (ICC = 0.87 to 0.99) (Table [3](#T3){ref-type="table"}). Reliability of accessibility ratings was less consistent across the categories. However, low cell sizes (n ≤ 5) for most of these analyses indicate that these results should be interpreted with some caution. ###### Reliability statistics for physical activity equipment separated by categories; Mean (SD) for participant data at Time 1 and Time 2, Intraclass Correlation Coefficient (ICC) and 95% Confidence Intervals (95% CI) or Mantel-Hanzel Chi Square test **Time 1** **Time 2** **T1 vs T2** ----------------------------------- ------------ -------------- -------------- ------------- --------------------------- **Sports Equipment**  \# of Items 23 34.6 (36.50) 23 33.5 (38.6) 0.96 (0.90 -- 0.98)  Accessibility   Away and difficult to get to 14.2 8.6 M-H χ^2^= 64.2, p \< 0.01   Away and easy to get to 25.2 27.7   In view and difficult to get to 3.5 5.1   In view and easy to get to 57.1 58.6 **Fitness Equipment**  \# of items 21 7.0 (4.69) 21 7.3 (5.93) 0.95 (0.88 -- 0.98)  Accessibility   Away and difficult to get to 6.6 2.2 M-H χ^2^= 5.63, p = 0.02   Away and easy to get to 28.6 31.1   In view and difficult to get to 3.3 2.2   In view and easy to get to 61.5 64.4 **Transportation Equipment**  \# of items 23 7.3 (4.58) 22 7.8 (5.42) 0.88 (0.74 -- 0.94)  Accessibility   Away and difficult to get to 2.9 0.0 M-H χ^2^= 20.4, p \< 0.01   Away and easy to get to 25.7 28.2   In view and difficult to get to 4.3 2.8   In view and easy to get to 67.1 69.0 **Water Sports Equipment**  \# of items 15 5.5 (6.09) 15 5.8 (9.01) 0.99 (0.97 -- 1.00)  Accessibility   Away and difficult to get to 12.5 7.7 M-H χ^2^= 0.02, p = 0.87   Away and easy to get to 45.8 26.9   In view and difficult to get to 4.2 3.9   In view and easy to get to 37.5 61.6 **Outdoor/Yard Equipment**  \# of items 20 12.5 (9.90) 20 13.5 (9.46) 0.87 (0.70 -- 0.95)  Accessibility   Away and difficult to get to 13.7 6.1 M-H χ^2^= 8.3, p \< 0.01   Away and easy to get to 9.2 20.6   In view and difficult to get to 4.6 2.3   In view and easy to get to 72.5 71.0 The test-retest reliability for the additional media questions ranged from K = 0.42 (0.10 to 0.73) for the number of video and computer games to K = 1.00 for the type of television and for the type of internet service. The Kappa statistics for the number of videos/DVDs (K = 0.60, 0.36 to 0.84) and number of television channels (K = 0.87, 0.71 to 1.00) were both acceptable. The Kappa statistic for television size (K = 0.71, 0.35 to 1.00) is based on categorizing the responses into small (\<25\"), medium (25--39\"), and large (\> = 40\"). Validity -------- Table [4](#T4){ref-type="table"} contains the results of the validation analyses. Associations between the participant and research assistants were moderate to strong across all variables (r = 0.67 to 0.98) and there were no significant differences between the mean values recorded by participants and research assistants (t-test p-values 0.20 to 0.72). For the accessibility ratings, the participants, compared to the research assistants, recorded a greater percentage of physical activity items as \"put away and difficult to get to\" and a smaller percentage of items as \"in plain view and easy to get to\" (p \< 0.001). No differences were observed for the accessibility ratings of the media items (p = 0.23). The PAASS was significantly greater for the research assistants compared to the participants (p = 0.03). There was no mean difference in the MAASS (p = 0.30). ###### Validity of PAMI variables; Mean ± SD; (N = 31) mean ± SD with Pearson Correlation Coefficient and p-value for independent t-test or percent with Mantel-Hanzel Chi Square **Variable** **Participant Time 1** **Research Assistant** **Correlation (Pearson r)** **t-test p-value** ----------------------------------- ------------------------ ------------------------ ----------------------------- -------------------- Total number of rooms 12.4 ± 2.91 12.0 ± 2.42 0.72\* 0.34 Physical Activity Equipment  \# of items 65.2 ± 54.12 74.8 ± 64.77 0.98\* 0.53  Household Density 5.0 ± 3.95 6.2 ± 5.54 0.93\* 0.33  Ave \# of Items/Bedroom 2.1 ± 2.68 2.8 ± 4.04 0.67\* 0.25  Checklist Quantity 16.8 ± 7.46 17.2 ± 7.27 0.94\* 0.36  Accessibility   Away and difficult to get to 9.5 2.5 \-- M-H χ^2^= 16.5,   Away and easy to get to 20.1 22 \-- p \< 0.001   In view and difficult to get to 5.3 2 \--   In view and easy to get to 65 73.6 \--  PAASS 209.9+162.24 261.2+231.3 0.93\* 0.03 Media Equipment  \# of items 8.5 ± 4.49 8.8 ± 4.47 0.93\* 0.20  Household Density 0.71 ± 0.34 0.75 ± 0.34 0.79\* 0.33  Ave \# of Items/Bedroom 1.7 ± 1.20 1.7 ± 1.20 0.94\* 0.49  \# Televisions in Home 2.5+1.23 2.9+1.44 0.87\* 0.25  \# Televisions in Bedrooms 1.16+1.04 1.26+1.09 0.96\* 0.72  Accessibility   Away and difficult to get to 1.0 1.3 \-- M-H χ^2^= 1.4,   Away and easy to get to 9.3 6.5 \-- p = 0.23   In view and difficult to get to 2.5 0.0 \--   In view and easy to get to 87.3 92.2 \--  MAASS 31.9+25.53 34.0+17.27 0.88\* 0.30 PAASS, Physical activity Availability and Accessibility Summary Score MAASS, Media Availability and Accessibility Summary Score \* p \< = 0.05 Table [5](#T5){ref-type="table"} contains the results comparing the participant (at time 1) and the research assistant separated by the physical activity equipment categories. Across categories, associations were generally high (r = 0.80 to 0.98) except for the \"transportation equipment\" category (r = 0.68). Research assistants recorded more sports equipment (t-test p \< 0.01), compared to the participants. Comparisons of the accessibility ratings indicate that accessibility for fitness, transportation, and water sports equipment were rated similarly between participants and research assistants (p = 0.23 to 0.73). However, research assistants recorded a greater proportion of sports and outdoor/yard equipment items as \"in view and easy to get to\" and a smaller proportion of items as \"put away and difficult to get to\". ###### Validity statistics for physical activity equipment separated by categories; Mean (SD) for participant at Time 1 and Research Assistants with Pearson correlation coefficients, t-test p-value or Mantel-Hanzel Chi Square test **Time 1** **RA** **T1 vs RA** ----------------------------------- ------------ -------------- -------------- -------------- ---------------------------- --------- **Sports Equipment**  \# of Items 31 33.5 (34.20) 31 42.9 (45.53) 0.98 \< 0.01  Accessibility   Away and difficult to get to 10.4 2.4 M-H χ^2^= 64.2, p \< 0.01   Away and easy to get to 24.5 28.0   In view and difficult to get to 3.7 0.8   In view and easy to get to 61.4 68.8 **Fitness Equipment**  \# of items 28 6.4 (4.31) 28 6.4 (4.78) 0.82 0.99  Accessibility   Away and difficult to get to 5.9 4.8 M-H χ^2^= 0.11, p = 0.73   Away and easy to get to 25.5 26.7   In view and difficult to get to 3.9 0.0   In view and easy to get to 64.7 68.6 **Transportation Equipment**  \# of items 28 7.7 (5.44) 31 7.8 (4.71) 0.68 0.90  Accessibility   Away and difficult to get to 4.3 2.2 M-H χ^2^= 1.41, p = 0.23   Away and easy to get to 16.1 14.3   In view and difficult to get to 10.8 5.5   In view and easy to get to 68.8 78.0 **Water Sports Equipment**  \# of items 16 4.7 (5.75) 17 4.6 (5.58) 0.92 0.75  Accessibility   Away and difficult to get to 12.0 0.0 M-H χ^2^= 0.29, p = 0.59   Away and easy to get to 44.0 52.0   In view and difficult to get to 4.0 8.0   In view and easy to get to 40.0 40.0 **Outdoor/Yard Equipment**  \# of items 28 12.8 (10.48) 28 14.4 (9.45) 0.80 0.17  Accessibility   Away and difficult to get to 14.3 2.8 M-H χ^2^= 12.63, p \< 0.01   Away and easy to get to 11.2 10.7   In view and difficult to get to 6.2 3.4   In view and easy to get to 68.3 83.1 Reliability and validity were also assessed for the overall home environment Activity:Media Ratio Score (the ratio of the PAASS to MAASS). At Time 1 and 2 the Activity:Media Ratio Score was 8.3 ± 7.98 and 8.3 ± 7.28, respectively (ICC = 0.91 (0.81 to 0.96)). There was a strong association between the participant- and research assistant-derived home environment scores (r = 0.94, P \< 0.01) and no significant difference in scores (participant; 8.2 ± 7.40, research assistant; 9.5 ± 10.44, p = 0.09). Discussion ========== This study assessed the test-retest reliability and validity of a self-administered, physical activity and media equipment inventory instrument for the home environment. Reliability ICCs for all of the physical activity and media equipment variables were moderate to high. This was true for all of the continuous variables, including the checklist quantity and the physical activity and media summary scores. When broken down by categories of physical activity items, reliability of the number of items remained consistently high across categories. The accessibility ratings were similar between the two administrations of the PAMI for fitness, transportation, and water sports equipment. In contrast, the accessibility ratings for sports and outdoor/yard equipment were different with fewer items considered \"put away and difficult to get to\" at the second administration, compared to the first. The reason for this difference may be a real change in where items in the home were placed or in the participants\' perception about the accessibility of the items. Validity, as assessed by comparing the participant and research assistant data, was consistently high for all calculated variables for both physical activity and media equipment. The PAASS summary score was higher for the research assistants compared to the participants and likely reflects the research assistants recording more items as \"in plain view and easy to get to\" compared to the participants. When broken down by categories of physical activity equipment, validity coefficients for the number of items in each category were still generally high. The distribution of accessibility ratings for the media equipment were similar between participants and research assistants with most items being \"in plain view and easy to get to\". Research assistants, however, recorded a smaller percentage of physical activity items as \"put away and difficult to get to\" and a greater percentage of items as \"in plain view and easy to get to\" as compared with participants. The difference between participant and research assistant ratings of accessibility across all physical activity items seems to be driven by a difference in ratings for the sports and yard/outdoor equipment. There may have been a discrepancy between what participants and research assistants considered \"put away\". While the participant may consider items placed on the floor by the kitchen door as, \"put away\" (if it was where those items are typically kept), the research assistants would have coded such items as, \"in plain view\". The term \"put away\" did not necessitate that the items be in a box, cabinet, drawer, or other closed container. In addition, this discrepancy may have occurred due to the protocol followed by the research assistants which allowed them to only record items that they could see, without opening doors, closets and looking under furniture (e.g., beds), unless the participant first did so. Our intention with implementing this protocol was to protect participant privacy. Further modification and testing of accessibility options is warranted given these findings. This was a small pilot study with several limitations. The majority of respondents were Caucasian, college educated and lived in single family homes. Additional validation studies are warranted in larger and more representative samples to determine the PAMI\'s psychometric properties in minority and lower income populations and those who live in apartments or condominiums. Participants were recruited from local parks and recreation departments which may have biased our sample by including relatively active families with relatively large amounts of physical activity equipment. Several families, however, had very few physical activity items, evidenced by the large standard deviations for most of the physical activity equipment variables. While the list of physical activity items was fairly comprehensive, some items may be more relevant in certain geographical areas (i.e., snow skis in Minnesota). Therefore, the list of items may need to be modified based on the population being studied. Only screen based media items were included on the PAMI. A more comprehensive assessment of media sources in the home would need to include magazines, radios, portable music players, cell phones, and possibly other sources. During the validation visits, participants walked from room to room to complete the PAMI, as the instructions indicate. Whether this happened during the second administration of the PAMI, when the participants were by themselves, is not known. The PAMI is also not designed to assess how frequently or who in the home, is using the equipment. However, in an ecological model, characterizing the environment is an important end in and of itself as the environment is important in influencing behavioral choices. Lastly, data were collected in the spring and seasonality should be considered when interpreting the location and accessibility data since some items (e.g., snow skis) may be moved and made very accessible during the winter but stored away at other times of the year. At present, it is not known if the PAMI is useful in characterizing homes along dimensions related to important health outcomes such as obesogeneity, levels of physical activity or sedentary behaviors. We also do not yet know if the PAMI is a useful instrument for showing change overtime as might be useful in intervention research that attempts to improve the healthfulness of the home environment. Future research will be needed to examine both the discriminant and predictive validity of the PAMI. In spite of these limitations, the PAMI is a potentially important addition to the field. The PAMI provides a more objective assessment of the presence of physical activity and media equipment items as compared to other research examining children\'s perception of adequate equipment \[[@B10],[@B11]\]. The PAMI instrument also expands on previous dichotomous checklist instruments \[[@B12],[@B13]\] by measuring the breadth (the number of unique items), magnitude (how many of each type of item) and ratings of accessibility of the items. Lessons learned --------------- In the PAMI instrument tested here (see Additional File [1](#S1){ref-type="supplementary-material"}), several types of foot wear were included as a category in the list of physical activity equipment (e.g., running shoes, hiking shoes, walking shoes, etc\...). There was considerable disagreement and confusion on identifying these particular types of shoes and, for many younger children and some adults, one pair of athletic shoes may serve multiple roles (walking, running, hiking). Also, there were so many pairs of shoes per household (range 1 to 59) that the discriminative ability of the number/accessibility of shoes was quite dubious. Therefore, types of shoes were not included in the analyses for this paper. Briefly, the ICC for \"shoes\" was 0.82 (0.62 -- 0.92) with a low correlation between the participant and the research assistant (r = 0.22). The participants reported slightly fewer pairs of shoes (14.7 ± 11.77), on average, compared to the research assistants (18.3 ± 12.25) (p = 0.07). For the main IDEA study, shoes were not included in the physical activity item list. In addition, based upon feedback from participants and research assistants, the PAMI equipment code numbers (included on a separate page of the tool) were sometimes difficult for the participants to find and record. Therefore, for the main IDEA study, the PAMI has been expanded so that each room has its own page containing the list of items where the participant indicates how many of each item are present in that room and the associated accessibility. We believe the re-formatted version of the PAMI will be easier for participants to complete, although this newer version was not tested in this pilot study. Conclusion ========== The PAMI may be a useful tool for describing the physical home environment related to opportunities for physical activity, sedentary behaviors and possibly the obesogeneity of the home environment. The main IDEA study will assess the associations between PAMI variables, physical activity and sedentary behaviors, and weight status of both the adult and youth participating in the study. In addition, the longitudinal design of the IDEA study will allow us to look at physical activity and media equipment in the home environment as potential determinants of physical activity, sedentary behavior, and as factors related to the obesogeneity of home environments. In conclusion, the results from this study indicate that the variables calculated for this study (number of items, density, checklist quantity, accessibility, summary scores) provide reliable estimates for describing the home environment related to the presence of physical activity and screen media equipment. The validity of the accessibility ratings for the physical activity items was less consistent and requires further investigation. Authors\' contributions ======================= JS and LL conceived the study. All authors (JS, MN, MP, LL) contributed significantly to the development of the inventory instrument and the study design, manuscript development and revision. All authors read and approved the final manuscript. Supplementary Material ====================== ###### Additional file 1 The Physical Activity and Media Inventory (PAMI) instrument. The full PAMI instrument used for this study. ###### Click here for file Acknowledgements ================ This research was supported by the University of Minnesota TREC (Transdisciplinary Research in Energetics and Cancer) Center Career Development Program (NCI Grant 1 U54 CA116849-01).

A Multiplication Based Logic Puzzle I really like this rhyme that I saw for the first time this week (even though it’s all over the net): Hey diddle diddle, the median’s the middle, You add then divide for the mean. The mode is the one that appears the most, And the range is the difference between. All of the numbers from 801 to 900 have at least 2 factors, but no more than 32 factors. 32 – 2 = 30, so 30 is the range of the amount of factors. There are 100 numbers from 801 to 900. If you list the amount of factors for each number, then arrange those amounts from smallest to largest, the amounts that will appear in the 50th and 51st spots will both be 6. That means that 6 is the median amount of factors. If we had different amounts in the 50th and 51st spots, we would average the two amounts together to get the median. If you add up the amounts of factors that the numbers from 801 to 900 have, you will get 794. If you divide 794 by 100, the number of entries, then you will know that 7.94 is the mean amount of factors. What about the mode? Which amount of factors appears the most? That’s why we are having a Horse Race, to see if more numbers have 2 factors, 3 factors, 4 factors, or a different amount of factors. So pick your pony. We’ll see which amount wins, and we’ll find out what the mode is at the same time. I should tell you that only perfect squares can have an odd amount of factors, so you probably don’t want to pick an odd amount. Here are some interesting facts about the numbers from 801 to 900 that might help you decide which pony to pick. We had the smallest two consecutive numbers with exactly 12 factors: (819, 820) We had the fourth prime decade: (821, 823, 827, 829). All four of those numbers are prime numbers and have exactly two factors. We had five consecutive numbers whose square roots can be reduced: (844, 845, 846, 847, 848). Three of those numbers had 6 factors, one had 10, and one had 12. We also had 840, the smallest number with exactly 32 factors 900 is the smallest number with exactly 27 factors. Coincidentally, the amount that is the mode will appear 27 times. As the following table shows, there are 42 integers from 801 to 900 that have square roots that can be simplified. 42 is more than any previous set of 100 numbers has given us. Even still we are still holding close to just under 40% of integers having square roots that can be simplified.

Uninterruptible power supplies (UPSs) including voltage converters are fundamental parts of many electrical systems such as power supply systems for computers and servers in data centers. UPSs can be used with many typical power systems including single and 3-phase connections, and can be used with low-power systems (e.g., a household computer) and high-power systems (e.g., large data centers or process facilities). High-power systems typically use a 3-phase power connection (e.g., X, Y, and Z phases). A 3-phase UPS voltage converter is typically used to provide 3-phase AC power to a 3-phase load, to convert a 3-phase AC voltage from one level to another, and to provide 3-phase power to a load in the event of a power failure. The input and output connections to a 3-phase UPS voltage converter are typically three- or four-terminal connections, one connection for each phase of the 3-phase power connection and an optional neutral connection. A battery is also typically coupled to the UPS voltage converter and is used to store energy for use in case of a power failure. Typical high power (e.g., above 100 kW) UPSs are operated using nominal AC input voltages of 3×400 V (in Europe) or 3×480 V (in the U.S.). Transformerless UPSs may operate with an internal DC bus voltage of ±450 V. In such a configuration, components contained in the UPS are preferably rated for at least 1200 V operation due to large voltage overshoots associated with stray inductances of physically large insulated gate bipolar transistor (IGBT) modules. The use of 1200 V components, however, typically leads to increased conduction and switching losses, thus lowering the efficiency.

How much of physical activity is recommended? WHO recommends that children and adolescents aged five to17 years: Should do at least 60 minutes of moderate to vigorous-intensity physical activity daily; physical activity of amounts greater than 60 minutes daily will provide additional health benefits; and should include activities that strengthen muscle and bone, at least three times per week. They recommend that adults aged 18 to 64 years: Should do at least 150 minutes of moderate-intensity physical activity throughout the week, or do at least 75 minutes of vigorous-intensity physical activity throughout the week, or an equivalent combination of moderate- and vigorous-intensity activity; for additional health benefits, adults should increase their moderate-intensity physical activity to 300 minutes per week, or equivalent; and muscle-strengthening activities should be done involving major muscle groups on 2 or more days a week. For adults aged 65 years and above, the WHO said they should do at least 150 minutes of moderate-intensity physical activity throughout the week, or at least 75 minutes of vigorous-intensity physical activity throughout the week, or an equivalent combination of moderate- and vigorous-intensity activity. For additional health benefits, the WHO said they should increase moderate intensity physical activity to 300 minutes per week, or equivalent; those with poor mobility should perform physical activity to enhance balance and prevent falls, three or more days per week; and muscle-strengthening activities should be done involving major muscle groups, two or more days a week. According to the WHO, the intensity of different forms of physical activity varies between people and in order to be beneficial for cardiorespiratory health, all activity should be performed in bouts of at least 10 minutes duration. The United Nations apex health body said regular physical activity of moderate intensity – such as walking, cycling, or doing sports – has significant benefits for health. “At all ages, the benefits of being physically active outweigh potential harm, for example through accidents. Some physical activity is better than doing none. By becoming more active throughout the day in relatively simple ways, people can quite easily achieve the recommended activity levels.” According to the WHO, regular and adequate levels of physical activity: improve muscular and cardiorespiratory fitness; improve bone and functional health; reduce the risk of hypertension, coronary heart disease, stroke, diabetes, breast and colon cancer and depression; reduce the risk of falls as well as hip or vertebral fractures; and are fundamental to energy balance and weight control. Culled From Vanguard. (Visited 30 times, 1 visits today) JOIN THE NEW DoctorsQuarters.com BBM CHANNEL ON C0015D291... Its New And Refreshingby by admin The Admin is a Medical Council Certified Medical Doctor trained in the prestigious Madonna University Elele, Nigeria. He is popular for being a fast rising online voice in Nigeria. I love to swim, read and play board games. I see myself as one who is destined to play a role in the way health services are rendered to the human race.

Q: How does the proxy mechanism work with proxy settings in browser We often find columns like Address, Port in web browser proxy settings. I know when we use proxy to visit a page, the web browser request the web page from the proxy server, but what I want to know is how the whole mechanism works? I have observed that many ISP allow only access to a single IP(of their website) after we exhausted our free data usage. But when we enter the site which we wants to browse in proxy URL and then type in the allowed IP, the site get loaded. How this works? A: In general, your browser simply connects to the proxy address & port instead of whatever IP address the DNS name resolved to. It then makes the web request as per normal. The web proxy reads the headers, uses the "Host" header of HTTP/1.1 to determine where the request is supposed to go, and then makes that request itself relaying all remaining data in both directions. Proxies will typically also do caching so if another person requests the same page from that proxy, it can just return the previous result. (This is simplified -- caching is a complex topic.) Since the proxy is in complete control of the connection, it can choose to route the request elsewhere, scrape request and reply data, inject other things (like ads), or block you altogether. Use SSL to protect against this. Some web proxies are "transparent". They reside on a gateway through which all IP traffic must pass and use the machine's networking stack to redirect outgoing connections to port 80 to a local port instead. It then behaves the same as though a proxy was defined in the browser. Other proxies, like SOCKS, have a dedicated protocol that allows non-HTTP requests to be made as well.

Summary: A new study reveals a diverse array of genetic changes that occur in the brain following sensory experiences. Source: Harvard. “Nature and nurture is a convenient jingle of words, for it separates under two distinct heads the innumerable elements of which personality is composed. Nature is all that a man brings with himself into the world; nurture is every influence from without that affects him after his birth.” – Francis Galton, cousin of Charles Darwin, 1874. Is it nature or nurture that ultimately shapes a human? Are actions and behaviors a result of genes or environment? Variations of these questions have been explored by countless philosophers and scientists across millennia. Yet, as biologists continue to better understand the mechanisms that underlie brain function, it is increasingly apparent that this long-debated dichotomy may be no dichotomy at all. In a study published in Nature Neuroscience on Jan. 21, neuroscientists and systems biologists from Harvard Medical School reveal just how inexorably interwoven nature and nurture are in the mouse brain. Using novel technologies developed at HMS, the team looked at how a single sensory experience affects gene expression in the brain by analyzing more than 114,000 individual cells in the mouse visual cortex before and after exposure to light. Their findings revealed a dramatic and diverse landscape of gene expression changes across all cell types, involving 611 different genes, many linked to neural connectivity and the brain’s ability to rewire itself to learn and adapt. The results offer insights into how bursts of neuronal activity that last only milliseconds trigger lasting changes in the brain, and open new fields of exploration for efforts to understand how the brain works. “What we found is, in a sense, amazing. In response to visual stimulation, virtually every cell in the visual cortex is responding in a different way,” said co-senior author Michael Greenberg, the Nathan Marsh Pusey Professor of Neurobiology and chair of the Department of Neurobiology at HMS. “This in essence addresses the long-asked question about nature and nurture: Is it genes or environment? It’s both, and this is how they come together,” he said. One out of many Neuroscientists have known that stimuli–sensory experiences such as touch or sound, metabolic changes, injury and other environmental experiences–can trigger the activation of genetic programs within the brain. Composed of a vast array of different cells, the brain depends on a complex orchestra of cellular functions to carry out its tasks. Scientists have long sought to understand how individual cells respond to various stimuli. However, due to technological limitations, previous genetic studies largely focused on mixed populations of cells, obscuring critical nuances in cellular behavior. To build a more comprehensive picture, Greenberg teamed with co-corresponding author Bernardo Sabatini, the Alice and Rodman W. Moorhead III Professor of Neurobiology at HMS, and Allon Klein, assistant professor of systems biology at HMS. Spearheaded by co-lead authors Sinisa Hrvatin, a postdoctoral fellow in the Greenberg lab, Daniel Hochbaum, a postdoctoral fellow in the Sabatini lab and M. Aurel Nagy, an MD-PhD student in the Greenberg lab, the researchers first housed mice in complete darkness to quiet the visual cortex, the area of the brain that controls vision. They then exposed the mice to light and studied how it affected genes within the brain. Using technology developed by the Klein lab known as inDrops, they tracked which genes got turned on or off in tens of thousands of individual cells before and after light exposure. The team found significant changes in gene expression after light exposure in all cell types in the visual cortex–both neurons and, unexpectedly, nonneuronal cells such as astrocytes, macrophages and muscle cells that line blood vessels in the brain. Roughly 50 to 70 percent of excitatory neurons, for example, exhibited changes regardless of their location or function. Remarkably, the authors said, a large proportion of non-neuronal cells–almost half of all astrocytes, for example–also exhibited changes. The team identified thousands of genes with altered expression patterns after light exposure, and 611 genes that had at least two-fold increases or decreases. Many of these genes have been previously linked to structural remodeling in the brain, suggesting that virtually the entire visual cortex, including the vasculature and muscle cell types, may undergo genetically controlled rewiring in response to a sensory experience. There has been some controversy among neuroscientists over whether gene expression could functionally control plasticity or connectivity between neurons. “I think our study strongly suggests that this is the case, and that each cell has a unique genetic program that’s tailored to the function of a given cell within a neural circuit,” Greenberg said. Question goldmine These findings open a wide range of avenues for further study, the authors said. For example, how genetic programs affect the function of specific cell types, how they vary early or later in life and how dysfunction in these programs might contribute to disease, all of which could help scientists learn more about the fundamental workings of the brain. “Experience and environmental stimuli appear to almost constantly affect gene expression and function throughout the brain. This may help us to understand how processes such as learning and memory formation, which require long-term changes in the brain, arise from the short bursts of electrical activity through which neurons signal to each other,” Greenberg said. One especially interesting area of inquiry, according to Greenberg, includes the regulatory elements that control the expression of genes in response to sensory experience. In a paper published earlier this year in Molecular Cell, he and his team explored the activity of the FOS/JUN protein complex, which is expressed across many different cell types in the brain but appears to regulate unique programs in each different cell type. Identifying the regulatory elements that control gene expression is critical because they may account for differences in brain function from one human to another, and may also underlie disorders such as autism, schizophrenia and bipolar disease, the researchers said. “We’re sitting on a goldmine of questions that can help us better understand how the brain works,” Greenberg said. “And there is a whole field of exploration waiting to be tapped.” About this neuroscience research article Additional authors on the study include Marcelo Cicconet, Keiramarie Robertson, Lucas Cheadle, Rapolas Zilionis, Alex Ratner and Rebeca Borges-Monroy. Funding: This work was supported by the National Institutes of Health (R01NS028829, R01NS046579, T32GM007753, R33CA212697, 5T32AG000222-23), F. Hoffmann-La Roche Ltd., the William F. Milton Fund, a Burroughs Wellcome Fund Career Award and an Edward J. Mallinckrodt Scholarship. Source: Ekaterina Pesheva – Harvard Publisher: Organized by NeuroscienceNews.com. Image Source: NeuroscienceNews.com image is credited to Lichtman Lab, Harvard University. Original Research: Abstract in Nature Neuroscience. DOI:10.1038/s41593-017-0029-5 Cite This NeuroscienceNews.com Article [cbtabs][cbtab title=”MLA”]Harvard “Nature, Meet Nurture.” NeuroscienceNews. NeuroscienceNews, 8 February 2018. <https://neurosciencenews.com/genetic-nature-experience-8455/>.[/cbtab][cbtab title=”APA”]Harvard (2018, February 8). Nature, Meet Nurture. NeuroscienceNews. Retrieved February 8, 2018 from https://neurosciencenews.com/genetic-nature-experience-8455/[/cbtab][cbtab title=”Chicago”]Harvard “Nature, Meet Nurture.” https://neurosciencenews.com/genetic-nature-experience-8455/ (accessed February 8, 2018).[/cbtab][/cbtabs] Abstract Single-cell analysis of experience-dependent transcriptomic states in the mouse visual cortex Activity-dependent transcriptional responses shape cortical function. However, a comprehensive understanding of the diversity of these responses across the full range of cortical cell types, and how these changes contribute to neuronal plasticity and disease, is lacking. To investigate the breadth of transcriptional changes that occur across cell types in the mouse visual cortex after exposure to light, we applied high-throughput single-cell RNA sequencing. We identified significant and divergent transcriptional responses to stimulation in each of the 30 cell types characterized, thus revealing 611 stimulus-responsive genes. Excitatory pyramidal neurons exhibited inter- and intralaminar heterogeneity in the induction of stimulus-responsive genes. Non-neuronal cells showed clear transcriptional responses that may regulate experience-dependent changes in neurovascular coupling and myelination. Together, these results reveal the dynamic landscape of the stimulus-dependent transcriptional changes occurring across cell types in the visual cortex; these changes are probably critical for cortical function and may be sites of deregulation in developmental brain disorders. Feel free to share this Neuroscience News.

NSSH Part 618 (Subpart A) Soil Properties and Qualities Subpart A – General Information 618.0 Definition and Purpose Soil properties are measured or inferred from direct observations in the field or laboratory. Examples of soil properties are particle-size distribution, cation-exchange capacity, and salinity. Soil qualities are behavior and performance attributes that are not directly measured. They are inferred from observations of dynamic conditions and from soil properties. Examples of soil qualities are corrosivity, natural drainage, frost action, and wind erodibility. Soil properties and soil qualities are the criteria used in soil interpretations, as predictors of soil behavior, and for classification and mapping of soils. The soil properties entered in the National Soil Information System (NASIS) must be representative of the soil and the dominant land use for which the interpretations are based. 618.1 Policy and Responsibilities Soil property data are collected, tested, and correlated as part of soil survey operations. These data are reviewed, supplemented, and revised as necessary. The Major Land Resource Area soil survey office (MLRA SSO) is responsible for collecting, testing, and correlating soil property data and interpretive criteria. The MLRA soil survey regional office (MO) is responsible for the development, maintenance, quality assurance, correlation, and coordination of the collection of soil property data that are used as interpretive criteria. This includes all the data elements listed below. The National Soil Survey Center (NSSC) is responsible for the training, review, and periodic update of soil interpretation technologies. The state soil scientist is responsible for ensuring that the soil interpretations are adequate for the Field Office Technical Guide and that they meet the needs of Federal, State, and local programs. 618.2 Collecting, Testing, and Populating Soil Property Data The collection and testing of soil property data is based on the needs described in the soil survey memorandum of understanding for individual soil survey areas. The collection and testing must conform to the procedures and guides established in this handbook. As aggregated component data, soil properties and qualities that are populated in NASIS are not meant to be site-specific. They represent the component as it occurs throughout the extent of the map unit. Most data entries are developed by aggregating information from point data (pedons) to create low, high, and representative values for the component. 618.3 Soil Properties and Soil Qualities The following sections list soil properties and qualities in alphabetical order and provide some grouping for climatic and engineering properties and classes. A definition, classes, significance, method, and guidance for NASIS database entry are given. The listing includes the soil properties and qualities in the NASIS database. For more details on the NASIS database, refer to part 639 of this handbook. For specifics on data structure, attributes, and choices in NASIS, refer to the NASIS webpage. Previous databases of soil survey information used metric or English units for soil properties and qualities. Values in English units, except for crop yields, were converted into metric units during transfer into the NASIS database. All future edits and entries in NASIS except yields and acreage will use metric units. Ranges of soil properties and qualities that are posted in the NASIS database for map unit components may extend beyond the established limits of the taxon from which the component gets its name, but only to the extent that interpretations do not change. However, the representative value (rv) is within the range of the taxon. 618.4 Albedo, Dry Definition.—“Albedo, dry” is the estimated ratio of the incident shortwave (solar) radiation that is reflected by the air-dry, less than 2 mm fraction of the soil surface to that received by it. Significance Soil albedo, as a function of soil color and angle of incidence of the solar radiation, depends on the inherent color of the parent material, organic matter content, and weathering conditions. Estimates of the evapotranspiration rates and predictions of soil water balances require albedo values. Evapotranspiration and soil hydrology models that are part of water quality and resource assessment programs require this information. Measurement.—There are instruments that measure albedo. Estimates.—Approximate the values by use of the following formula: Soil Albedo = 0.069 × (Color Value) - 0.114. For albedo, dry, use dry color value. Surface roughness has a separate significant impact on the actual albedo. The equation above is the albedo of <2.0 mm smoothed soil condition; if the surface is rough because of tillage, the albedo differs. Entries.—Enter the high, low, and representative values of albedo for the map unit component. The range of valid entries is from 0 to 1, and hundredths (two decimal places) are allowed. 618.5 Artifacts in the Soil Definition.—“Artifacts” are objects or materials created or modified by humans, usually for a practical purpose in habitation, manufacturing, excavation, or construction activities. Examples of artifacts include bitumen (asphalt), brick, concrete, metal, paper, plastic, rubber, and wood products. Artifacts are commonly referred to as “discrete artifacts” if they are 2 mm or larger in diameter and are not compacted into a root-limiting layer that impedes root growth or water movement. Significance.—Artifacts can constitute a significant portion of the soil. The amount and type of particulate artifacts can contribute substantially to various trace metals and total carbon contents of soils. Discrete artifacts which are both cohesive and persistent, defined below, are treated in a similar manner as rock fragments when populating the standard sieves or in calculations involving sieve entries. Discrete artifacts which are noncohesive, nonpersistent, or both are not considered fragments for sieve entries or calculations involving those entries. Measurement.—The fraction from 2 to 75 mm in diameter may be measured in the field. However, 50 to 60 kg of sample material may be necessary if there is an appreciable amount of fragments near 75 mm. An alternative means of measuring is to visually estimate the volume of the 20 to 75 mm fraction, then sieve and weigh the 2 to 20 mm fraction. The fraction 75 mm (3 inches) or greater is usually not included in soil samples taken in the field for laboratory testing. Measurements can be made in the field by weighing the dry sample and the portion retained on a 3-inch screen. The smallest dimension of discrete artifacts is used to determine whether these items pass through a sieve. The quantity is expressed as a weight percentage of the total soil. A sample as large as 200 pounds to more than a ton may be needed to assure that the results are representative. Measurements of the fraction from 75 to 250 mm (3 to 10 inches) and the fraction greater than 250 mm (10 inches) in diameter are usually obtained from volume estimates. Estimates Estimates of discrete artifacts are made similarly to the way estimates of rock fragments are made. These estimates are usually made by visual means and are on the basis of percent by volume. The percent by volume is converted to percent by weight by using the average bulk unit weights for the soil and the specific artifacts. These estimates are made during investigation and mapping activities in the field. They are expressed as ranges that include the estimating accuracy as well as the range of values for a component. Treated and untreated wood products (e.g., lumber) are considered artifacts. They are not considered wood fragments such as those associated with the woody materials (e.g., tree branches) described in organic soils. Measurements or estimates of discrete artifacts less than strongly cemented are made prior to any rolling or crushing of the sample. Artifact Cohesion Definition.—“Artifact Cohesion” is the relative ability of the artifact to remain intact after significant disturbance. Significance.—Artifacts that break down easily are similar to pararock fragments in that these artifacts break down to become part of the fine-earth fraction of the soil. Noncohesive artifacts are excluded from entries for the standards sieves and are not used in sieve calculations. Entries.—Enter cohesive or noncohesive in the Component Horizon Human Artifacts and the Pedon Horizon Human Artifacts tables of the NASIS database. Cohesion is based on whether the artifact can be easily broken into <2 mm size pieces either in the hands or with a mortar and pestle. Artifacts that cannot easily be broken are cohesive. All others are considered noncohesive. Artifact Kind Definition.—“Artifact Kind” is the type of object or material being described. Significance.—Each type of artifact is associated with a combination of other property entries that is used to determine whether the artifact is considered for sieve entries and calculations. The type of artifact also gives clues to the age of the deposit as well as the potential toxicity. Entries.—Enter the artifact kind in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables. Enter the appropriate choice for the kind of discrete artifact from the following list— Bitumen (asphalt) Boiler slag Bottom ash Brick Cardboard Carpet Cloth Coal combustion by-products Concrete Debitage Fly ash Glass Metal Paper Plasterboard Plastic Potsherd Rubber Treated wood Untreated wood. Artifact Penetrability Definition.—“Artifact Penetrability” is the relative ease with which roots can penetrate the artifact and potentially extract any stored moisture, nutrients, or toxic elements. Significance.—Artifacts that are penetrable may increase the available water holding capacity of a soil and should be factored in such calculations. The availability of supplemental nutrients and toxic elements is also greatest in penetrable artifacts. Entries.—Enter nonpenetrable or penetrable in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables based on whether roots can penetrate the solid parts of the artifact or between the component parts of the artifact. Artifact Persistence Definition.—“Artifact Persistence” is the relative ability of solid artifacts to withstand weathering and decay over time. Significance.—Artifacts that decay quickly are similar to pararock fragments and are treated as such in sieve calculations. Entries.—Enter nonpersistent or persistent in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables based on whether the artifact is expected to decay in less than a decade or greater than a decade. Nonpersistent artifacts are expected to decay in less than a decade. Persistent artifacts remain intact for a decade or more. Artifact Roundness Definition.—“Artifact Roundness” is an expression of the sharpness of edges and corners of objects. Significance.— The roundness of artifacts impacts water infiltration, root penetration, and macropore space. Classes.—The artifact roundness classes follow those used for fragment roundness: Roundness Class Definition Very angular Strongly developed faces with very sharp, broken edges. Angular Strongly developed faces with sharp edges (Soil Survey Manual (SSM)). Subangular Detectable flat faces with slightly rounded corners. Subrounded Detectable flat faces with well-rounded corners (SSM). Rounded Flat faces absent or nearly absent with all corners rounded (SSM). Well rounded Flat faces absent with all corners rounded. Entries.—Enter the appropriate artifact roundness class name for the record of artifacts populated in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables. Artifact Safety Definition.—“Artifact Safety” is the degree of risk to humans from contact with soils that contain artifacts. Physical contact with soils containing dangerous or harmful artifacts should be avoided unless proper training and protective clothing is available. The risk is based on toxicity to living organisms and not the physical risk that may be present from sharp or heavy objects. Harmful toxicity may be immediate or long-term, or through direct or indirect contact. Examples of innocuous artifacts include brick, concrete, glass, plastic, unprinted paper and cardboard, and untreated wood. Some examples of noxious artifacts are batteries, bitumen (asphalt), fly ash, garbage, paper printed with metallic ink, and wood treated with arsenic. Significance.—Noxious artifacts are dangerous and require special handling when sampling. Areas with noxious artifacts should have restricted human contact. Entries.—Enter innocuous or noxious in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables based on whether the artifacts are potentially toxic to living beings. Artifact Shape Definition.—“Artifact Shape” is a description of the overall shape of the object. Significance.—Artifact shape differs from rock, pararock, and wood fragment shape descriptions and is important for fluid flow in the soil as well as influencing excavation difficulty. Entries.—Enter the appropriate artifact shape class name for each record of artifacts populated in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables. Artifact Size Definition.—“Artifact Size” is based on the cross-sectional diameter of the object. Significance.—The size of discrete artifacts is significant to the use and management of the soil. Artifact sizes ranging from 2 mm to 75 mm that are both cohesive and persistent are considered when estimating the percent passing the sieves. It affects equipment use, excavation, construction, and recreational uses. Entries.—Enter the cross-sectional diameter size of the ≥ 2 mm artifacts described in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables. The range of valid entries is from 2 to 3,000 millimeters, and only whole numbers (integers) are allowed. Artifact Volume Definition.—“Artifact Volume” is the volume percentage of the horizon occupied by the 2 mm or larger fraction (20 mm or larger for wood artifacts) on a whole soil base. Entries.—Enter the high, low, and representative values for the percent volume present of each size class and kind of artifact populated in the Component Horizon Human Artifacts and Pedon Horizon Human Artifacts tables. The range of valid entries is from 0 to 100 percent, and only whole numbers (integers) are allowed. 618.6 Available Water Capacity Definition.—“Available water capacity” is the volume of water that should be available to plants if the soil, inclusive of fragments, were at field capacity. It is commonly estimated as the amount of water held between field capacity and wilting point, with corrections for salinity, fragments, and rooting depth. Classes.—Classes of available water capacity are not normally used except as adjective ratings that reflect the sum of available water capacity in inches to some arbitrary depth. Class limits vary according to climate zones and the crops commonly grown in the areas. The depth of measurement also is variable. Significance.—Available water capacity is an important soil property in developing water budgets, predicting droughtiness, designing and operating irrigation systems, designing drainage systems, protecting water resources, and predicting yields. Estimates.—The most common estimates of available water capacity are made in the field or the laboratory as follows: Field capacity is determined by sampling the soil moisture content just after the soil has drained following a period of rain and humid weather, after a spring thaw, or after heavy irrigation. Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS, provides more information. The 15-bar moisture content of the samples is determined with pressure membrane apparatus. An approximation of soil moisture content at field capacity is commonly made in the laboratory using 1/3-bar moisture percentage for clayey and loamy soil materials and 1/10-bar for sandy materials. Recently, some soil physicists have been using 1/10-bar instead of 1/3-bar for clayey and loamy soil materials and 1/20-bar for sandy soil materials. If data are available, estimates are based on available water capacity measurements. If data are not available, data from similar soils are used as a guide. The relationship between available water capacity and other soil properties has been studied by many researchers. Soil properties that influence available water capacity are particle size; size, shape, and distribution of pores; organic matter; type of clay mineral; and structure. If roots are excluded from a horizon such as a duripan, the amount of water available to plants is nearly zero. Available water capacity values are zero for layers that exclude roots. If roots are restricted but not excluded, estimates of available water capacity are reduced according to the amount of dense material in the layers and the space available for root penetration. Depending on the ability of roots to enter the soil mass and utilize the water, values for the soils with these dense layers may be significantly less than for soils of similar texture that do not have pans. Entries are made for all soil layers below dense layers only if roots are present. Depending on their abundance and porosity, rock and pararock fragments reduce available water capacity. Nonporous fragments reduce available water capacity in proportion to the volume they occupy. For example, 50-percent nonporous cobbles reduce available water capacity as much as 50 percent. Porous fragments, such as sandstone, may reduce available water capacity to a lesser extent. Several factors contribute to a lower amount of plant growth on saline soils. However, as a rough guide, available water capacity is reduced by about 25 percent per 4 mmhos cm-1 electrolytic conductivity of the saturated extract. Soils high in gibbsite or kaolinite, such as Oxisols and Ultisols, may have available water capacity values that are about 20 percent lower than those with equal amounts of 2:1 layer-lattice clays. Soils high in organic matter have a higher available water capacity than soils low in organic matter if the other properties are the same. Entries.—Enter high, low, and representative values for available water capacity in cm per cm for each horizon. Enter “0” for layers that exclude roots. The range of valid entries is from 0 to 0.7 cm per cm, and hundredths (two decimal places) are allowed. 618.7 Bulk Density, One-Third Bar Definition.—“Bulk density, one-third bar” is the oven-dried weight of the less than 2 mm soil material per unit volume of soil at a water tension of 1/3 bar (33 kPa). Significance.—Bulk density influences plant growth and engineering applications. It is used to convert measurements from a weight basis to a volume basis. Within a family level particle-size class, bulk density is an indicator of how well plant roots are able to extend into the soil. Bulk density is used to calculate porosity. Bulk density at 33 kPa is used for soil classification in the required characteristics for andic soil properties and in the criteria for Andic, Aquandic, and Vitrandic subgroups. Plant Growth.—Bulk density is an indicator of how well plant roots are able to extend into the soil. Root restriction initiation and root-limiting bulk densities are shown below for various particle-size classes. Engineering Applications.—Soil horizons with bulk densities less than those indicated below have low strength and would be subject to collapse if wetted to field capacity or above without loading. They may require special designs for certain foundations. Particle-Size Class Bulk Density (g cm-3) Sandy <1.60 Loamy coarse-loamy <1.40 fine-loamy <1.40 coarse-silty <1.30 fine-silty <1.40 Clayey <1.10 Estimates.—The weight applies to the oven-dry soil, and the volume applies to the soil at or near field capacity. Bulk density is a use-dependent property. The entry should represent the dominant use for the soil. Entries.—Enter bulk density at one-third bar with the low, high, and representative values for each horizon. The range of valid entries is from 0.02 to 2.6 g cm-3, and hundredths (two decimal places) are allowed. Values should be estimated to the nearest 0.05 g cm-3. 618.8 Bulk Density, Oven Dry Definition.—“Bulk density, oven dry” (pbod) is the oven-dry weight of the less than 2 mm soil material per unit volume of oven-dry soil. Estimates.—The value pbod is derived by the following formula: pbod = [(linear extensibility percent/100) + 1]3 Where, linear extensibility percent is adjusted to a <2 mm basis. Entries.—Enter the high, low, and representative values for each horizon. The range of valid entries is from 0.02 to 2.6 g cm-3, and hundredths (two decimal places) are allowed. 618.9 Bulk Density, Satiated Definition.—“Bulk density, satiated” (pbsat) is the oven-dry weight of the less than 2 mm soil material per unit volume of soil at a water tension of 0 bar. The measurement is only used for subaqueous soils. Significance.—Coastal wetland and subaqueous soils exist in their environment at saturation. Very low bulk density soils in submerged environments often contain a large percentage of water, making them very fluid. These soil qualities are important in subaqueous soil interpretations for shellfish and rooted vegetation habitat as well as constriction, dredge operations, and the calculation of carbon stocks. Estimates.—The value pbsat is calculated based on the dried weight of a known volume of soil at the field moisture status. Sampling methods can vary depending on environment. For samples taken as vibracores and opened by cutting, a 50 ml plastic syringe with the end removed is used to collect a mini-core. The plunger can be fixed at the 10 ml volume mark and the syringe gently pused into the split vibracore sample to collect a known volume of sample. This technique is a variation of the field-state core method. Entries.—Enter the high, low, and representative values for each horizon. The range of valid entries is from 0.02 to 2.6 g cm-3, and hundredths (two decimal places) are allowed. 618.10 Calcium Carbonate Equivalent Definition.—“Calcium carbonate equivalent” is the quantity of carbonate (CO3) in the soil expressed as CaCO3 and as a weight percentage of the less than 2 mm size fraction. Significance.—The availability of plant nutrients is influenced by the amount of carbonates in the soil. This is a result of the effect that carbonates have on soil pH and of the direct effect that carbonates have on nutrient availability. Nitrogen fertilizers should be incorporated into calcareous soils to prevent nitrite accumulation or ammonium-N volatilization. The availability of phosphorus and molybdenum is reduced by the high levels of calcium and magnesium which are associated with carbonates. In addition, iron, boron, zinc, and manganese deficiencies are common in soils that have a high calcium carbonate equivalent. In some climates, soils that have a high calcium carbonate equivalent in the surface layer are subject to wind erosion. This effect may occur in soils that have a calcium carbonate equivalent of more than 5 percent. A strongly or violently effervescent reaction to cold, dilute hydrochloric acid (HCL) defines calcareous in the wind erodibility groups because of the significance of finely divided carbonates. Calcium carbonate equivalent is used for soil classification in the criteria for several diagnostic horizons (e.g., mollic epipedon), Rendolls suborder, Rendollic Eutrudepts subgroup, and carbonatic mineralogy class. Measurement.—Calcium carbonate equivalent is measured by a method that uses an aqueous solution of 3 normal hydrogen chloride. The method is outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. It also may be measured in the field using calcimeters. Entries.—Enter the high, low, and representative values for each horizon listed. Round values to the nearest 5 percent for horizons that have more than 5 percent CaCO3 and to the nearest 1 percent for those with less than 5 percent. Enter 0 if the horizon does not have free carbonates. The range of valid entries is from 0 to 110 percent, and only whole numbers (integers) are allowed. 618.11 Cation-Exchange Capacity NH4OAc pH 7 Definition.—“Cation-exchange capacity” is the amount of exchangeable cations that a soil can adsorb at pH 7.0. Significance.—Cation-exchange capacity is a measure of the ability of a soil to retain cations, some of which are plant nutrients. Soils that have a low cation-exchange capacity hold fewer cations and may require more frequent applications of fertilizer than soils that have a high cation-exchange capacity. Soils that have a high cation-exchange capacity have the potential to retain cations, which reduces the risk of the pollution of ground water. Cation-exchange capacity is used indirectly in soil classification, when recalculated for just the noncarbonate clay fraction, in the required characteristics for kandic and oxic horizons and as a criterion for specific subgroup taxa in Alfisols (e.g., Kandic Paleustalfs), Entisols, Inceptisols, and Mollisols. Cation-exchange capacity is also used in calculating the ratio of cation-exchange capacity to percent noncarbonate clay for classifying certain soils at the family level into cation-exchange activity classes. The latest edition of the Keys to Soil Taxonomy has more information on applying this ratio in classification. Measurement.—Cation-exchange capacity is measured by the methods outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. The ammonium acetate method gives the cation-exchange capacity value (CEC-7) for soils that have pH >5.5 or contain soluble salts. This method uses a solution of one normal ammonium acetate buffered at pH 7.0 to provide the extracting index cation (NH4+). Cation-exchange capacity is reported, on a <2 mm base, in centimoles per kilogram (cmol(+) kg-1), which are equivalent to milliequivalents per 100 grams (meq 100 g-1) of fine-earth soil. If the pH is less than 5.5, use effective cation-exchange capacity (refer to Section 618.19). Entries.—Enter the high, low, and representative values of the estimated range in cation-exchange capacity, in meq 100 g-1, for each horizon with pH >5.5. The range of valid entries is from 0 to 400 meq 100 g-1, and tenths (one decimal place) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.99. 618.12 Climatic Setting Climatic setting includes frost-free period, precipitation, temperature, and evaporation. These elements are useful in determining the types of natural vegetation or crops that grow or can grow in an area and in planning management systems for vegetation. Climatic data are observed nationally by the National Weather Service Cooperative Network, which consists of approximately 10,000 climate stations. The records are available from the Climatic Data Access Facility (CDAF) at Portland, Oregon. Climatic data are delivered to the field through the Climatic Data Access Network. The Climatic Data Access Network consists of climatic data liaisons established in each State and at National Headquarters. Climatic data that are input into NASIS are obtained from the respective climatic data liaison. Climatic data may also be obtained from project weather stations or from the State climatologist. NRCS has selected the current standard “normal” period of 1971 to 2000 for climate database entries. Existing entries in the NASIS database may reflect the prior period of 1961 to 1990. Always check with your State’s climatic data liaison before using a climate station that has less than 30 years of records or that is located outside a county. Footnote the source of the data, the station, and the starting and ending year of record. Means are given as a range to represent the change of the climate over the geographic extent of the assigned soil. Frost-Free Period Definition.—“Frost-free period” is the expected number of days between the last freezing temperature (0 °C) in spring (January-July) and the first freezing temperature (0 °C) in fall (August-December). The number of days is based on the probability that the values for the standard “normal” period will be exceeded in 5 years out of 10. Entries.—Enter the high, low, and representative values for the map unit component. Enter 365 for each value for taxa that are frost-free all year and 0 for those that have no frost-free period. Entries are rounded to the nearest 5 days. Precipitation, Mean Annual Definition.—“Mean annual precipitation” is the arithmetic average of the total annual precipitation taken over the standard “normal” period. Precipitation refers to all forms of water, liquid or solid, that fall from the atmosphere and reach the ground. Entries.—Enter the high, low, and representative values in millimeters of water to represent the spatial range for the map unit component. The range of valid entries is from 0 to 11,500 mm, and only whole numbers (integers) are allowed. Air Temperature, Mean Annual Definition.—“Mean annual air temperature” is the arithmetic average of the daily maximum and minimum temperatures for a calendar year taken over the standard “normal” period. Entries.—Enter the high, low, and representative values for the map unit component to represent the spatial range in degrees Celsius (centigrade). The range of valid entries is from -50.0 to 50.0 degrees, and tenths (one decimal place) are allowed. Use a minus sign to indicate temperatures below zero. Daily Average Precipitation Definition.—“Daily average precipitation” is the total precipitation for the month divided by the number of days in the month for the standard “normal” period. Entries.—Enter the high, low, and representative values, in millimeters. The range of valid entries is from 0 to 750 mm. Record values to the nearest whole number (integer). Daily Average Potential Evapotranspiration Definition.—“Daily average potential evapotranspiration” is the total monthly potential evapotranspiration divided by the number of days in the month for the standard “normal” period. Entries.—Enter the high, low, and representative values for daily average potential evapotranspiration in millimeters. The range of valid entries is 0 to 300 mm. Record values to the nearest whole number (integer). 618.13 Corrosion Various metals and other materials corrode when they are on or in the soil, and some metals and materials corrode more rapidly when in contact with specific soils than when in contact with others. Corrosivity ratings are given for two of the common structural materials, uncoated steel and concrete. Uncoated steel Definition.—“Risk of corrosion for uncoated steel” is the susceptibility of uncoated steel to corrosion when in contact with the soil. Classes.—The classes for risk of corrosion to uncoated steel are low, moderate, and high. Significance.—Risk of corrosion to uncoated steel pertains to the potential soil-induced electrochemical or chemical action that converts iron into its ions, thereby dissolving or weakening uncoated steel. Guides.—Part 618, Subpart B, Exhibits, Section 618.80 gives the relationship of soil water, general texture group, acidity, and content of soluble salts (as indicated by either electrical resistivity at field capacity or electrolytic conductivity of the saturated extract of the soil) to corrosion classes. Soil reaction (pH) correlates poorly with corrosion potential; however, a pH of 4.0 or less almost always indicates a high corrosion potential. Ratings, which are based on a single soil property or quality, that place soils in relative classes for corrosion potential must be tempered by knowledge of other properties and qualities that affect corrosion. A study of soil properties in relation to local experiences with corrosion helps soil scientists and engineers to make soil interpretations. Special attention must be given to those soil properties that affect the access of oxygen and moisture to the metal, the electrolyte, the chemical reaction in the electrolyte, and the flow of current through the electrolyte. Special attention must be given to the presence of sulfides or of minerals, such as pyrite, that can be weathered readily and thus cause a high degree of corrosion in metals. The possibility of corrosion is greater for extensive installations that intersect soil boundaries or soil horizons than for installations that are in one kind of soil or in one soil horizon. Using interpretations for corrosion without considering the size of the metallic structure or the differential effects of using different metals may lead to wrong conclusions. Activities that alter the soil, such as construction, paving, fill and compaction, and surface additions, can increase the possibility of corrosion by creating an oxidation cell that accelerates corrosion. Mechanical agitation or excavation that results in aeration and in a discontinuous mixing of soil horizons may also increase the possibility of corrosion. Entries.—Enter the appropriate class of risk of corrosion for uncoated steel for the whole map unit component. The classes are low, moderate, and high. Concrete Definition.—“Risk of corrosion for concrete” is the susceptibility of concrete to corrosion when in contact with the soil. Classes.—The classes for risk of corrosion to concrete are low, moderate, and high Significance.—Risk of corrosion to concrete pertains to the potential soil-induced chemical reaction between a base (the concrete) and a weak acid (the soil solution). Special cements and methods of manufacturing may be used to reduce the rate of deterioration in soils that have a high risk of corrosion. The rate of deterioration depends on soil texture and acidity; the amount of sodium or magnesium sulfate present in the soil, singly or in combination; and the amount of sodium chloride (NaCl) present in the soil. The presence of NaCl is evaluated because it is used to identify the presence of seawater, rather than because of its corrosive effects on concrete. Seawater contains sulfates, which are one of the principal corrosive agents. A soil that has gypsum or other sulfate minerals requires a special cement in the concrete mix. The calcium ions in gypsum react with the cement and weaken the concrete. Entries.—Enter the appropriate class of risk of corrosion for concrete for the whole map unit component. The classes are low, moderate, and high. 618.14 Crop Name and Yield Definition.—“Crop name” is the common name for the crop. “Crop yield” is crop yield units per unit area for the specified crop. Classes.—The crop names and the units of measure for yields that are allowable as data entries are listed in the NASIS data dictionary. See Part 618, Subpart B, Exhibits, Section 618.82 for the Web address of the current NASIS data dictionary. Significance.—Crop names and units of measure are important as records of crop yield. Although the crops and yield often are specific to the time when the soil survey was completed, the ranking and comparison between soils within a soil survey are helpful. These crop and yield data are used to evaluate the soil productive capabilities, cash rent, and land values. Generally, only the most important crops are listed and only the best management is reflected. Estimates Crop names and yields are specific to the soil survey area. Although the listing of crop names is not limited to any number, only the most important crops in the survey area should be used. The yields are derived in a number of ways but should represent a high level of management by leading commercial farmers, which tends to produce the highest economic return per acre. This level of management includes using the best varieties; balancing plant populations and added plant nutrients to the potential of the soil; controlling erosion, weeds, insects, and diseases; maintaining optimum soil tilth; providing adequate soil drainage; and ensuring timely operations. Generally only a representative value is used for each map unit component for non-MLRA soil survey areas. MLRA soil survey areas use the high and low representative values from map unit components of non-MLRA soil survey areas. High and low values represent the range of representative values for a high level of management across the survey area or across several survey areas. Entries.—Enter the common crop name and units of measure. Enter the corresponding irrigated yields, nonirrigated yields, or both, as appropriate for the component. Yields can be posted as high, low, and representative values for the map unit component. 618.15 Diagnostic Horizon Feature — Depth to Bottom Definition.—The diagnostic horizon feature “depth to bottom” is the distance from the top of the soil to the base of the identified diagnostic horizon or to the lower limit of the occurrence of the diagnostic feature. Measurement.—Distance is measured from the top of the soil, which is defined as the top of the mineral soil, or, for soils with “O” horizons, the top of any organic layer that is at least slightly decomposed. For soils that are covered by 80 percent or more rock or pararock fragments, the top of the soil is the mean height of the top of the fragments. See chapter 3 in the Soil Survey Manual for a complete discussion. Entries.—The values for the diagnostic horizon feature depth to bottom used to populate component data in NASIS are not specific to any one point; they are a reflection of commonly observed values based on field observations and are intended to model the component as it occurs throughout the map unit. Enter the high, low, and representative values in whole centimeters. The high value represents either the greatest depth to which the base of the diagnostic horizon or feature extends or, for horizons for features extending beyond the limit of field observation, is the depth to which observation was made (usually no more than 200 cm). In the case of lithic contact, paralithic contact, and petroferric contact, the entries for depth to the bottom of the diagnostic feature will be the same as the entries for depth to the top of the feature since the contact has no thickness. 618.16 Diagnostic Horizon Feature — Depth to Top Definition.—The diagnostic horizon feature “depth to top” is the distance from the top of the soil to the upper boundary of the identified diagnostic horizon or to the upper limit of the occurrence of the diagnostic feature. Measurement.— Distance is measured from the top of the soil, which is defined as the top of the mineral soil, or, for soils with “O” horizons, the top of any organic layer that is at least slightly decomposed. For soils that are covered by 80 percent or more rock or pararock fragments, the top of the soil is the mean height of the top of the fragments. See chapter 3 in the Soil Survey Manual for a complete discussion. Entries.—The values for the diagnostic horizon feature depth to top used to populate component data in NASIS are not specific to any one point; they are a reflection of commonly observed values based on field observations and are intended to model the component as it occurs throughout the map unit. Enter the high, low, and representative values in whole centimeters. 618.17 Diagnostic Horizon Feature — Kind Definition.—The diagnostic horizon feature “kind” is the kind of diagnostic horizon or diagnostic feature present in the soil. Significance.—Diagnostic horizons and features are a particular set of observable or measurable soil properties, defined in Soil Taxonomy, that are used to classify a soil. They have been chosen because they are thought to be the marks left on the soil as a result of the dominant soil-forming processes. In many cases they are thought to occur in conjunction with other important accessory properties. The utilization of diagnostic horizons and features in the classification process allows the grouping of soils that have formed as a result of similar genetic processes. The grouping, however, is done on the basis of observable or measurable properties rather than by speculation about the genetic history of a particular soil. Entries.—The diagnostic horizons and features are listed in the latest edition of the Keys to Soil Taxonomy. Allowable terms are given in the NASIS data dictionary. 618.18 Drainage Class Definition.—“Drainage class” identifies the natural drainage condition of the soil. It refers to the frequency and duration of wet periods. Significance.—Drainage classes provide a guide to the limitations and potentials of the soil for field crops, forestry, range, wildlife, and recreational uses. The class roughly indicates the degree, frequency, and duration of wetness, which are factors in rating soils for various uses. Estimates.—Infer drainage classes from observations of landscape position and soil morphology. In many soils the depth and duration of wetness relate to the quantity, nature, and pattern of redoximorphic features. Correlate drainage classes and redoximorphic features through field observations of water tables, soil wetness, and landscape position. Record the drainage classes assigned to the series. Entries.—Enter the drainage class name for each map unit component. Use separate map unit components for different drainage class phases or for drained versus undrained phases where needed. Significance.—Cation-exchange capacity (CEC) is a measure of the ability of a soil to retain cations, some of which are plant nutrients. Soils that have a low cation-exchange capacity hold fewer cations and may require more frequent applications of fertilizer and amendments than soils that have a high cation-exchange capacity. Effective CEC (ECEC) is a measure of CEC that is particularly useful in soils whose ion-exchange capacity is largely a result of variable charge components, such as allophane, imogolite, kaolinite, halloysite, hydrous iron and aluminum oxides, and organic matter. As a result, the CEC of these soils is not a fixed number but is a function of pH. Examples of taxa commonly displaying pH-dependent charge include some Andisols, Histosols, acidic Inceptisols, Oxisols, Spodosols, and weathered Ultisols with kaolinitic or halloysitic mineralogies dominated by iron and aluminum oxyhydroxide minerals. Measurement.—Effective cation-exchange capacity is calculated from the analytical results of two separate laboratory methods. One method measures the basic cations (Ca2+, Mg2+, Na+, K+) extractable in a solution of one normal ammonium acetate buffered at pH 7.0. Another method measures the aluminum extractable in a solution of one normal potassium chloride(for soil horizons with a 1:1 water pH of 5.5 or less). The ECEC value is then calculated and reported for soil horizons that have pH <5.5 and that are low in soluble salts. For soils that have a pH of 5.5 or greater, the ECEC usually equals only the sum of the NH4OAc extractable bases. Manual ECEC population in NASIS for soil horizons with pH values between 5.6 and 7.0 is optional and is only needed if there is a significant difference from the populated CEC values (based on NH4OAc buffered at pH 7.0). An alternate procedure exists to measure ECEC. It involves a direct measurement by using a neutral unbuffered salt (NH4Cl) and is an analytically determined value. For a soil with a pH of less than 7.0 (in water, 1:1), the ECEC value should be less than the CEC value measured with a buffered solution at pH 7.0. The ECEC by NH4Cl is equal to the NH4OAc extractable bases plus the KCl extractable Al for noncalcareous soils. For more discussion on ECEC see Soil Survey Investigations Report No. 45, Soil Survey Laboratory Information Manual, Version 2.0, February 2011, USDA, NRCS. The laboratory methods for both the standard and alternate procedures are outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Effective cation-exchange capacity is reported, on a <2 mm base, in centimoles per kilogram (cmol(+) kg-1) of soil, which are equivalent to milliequivalents per 100 grams (meq 100 g-1) of fine-earth soil. Entries.—Enter the high, low, and representative values of the estimated range in effective cation-exchange capacity at the field pH of the soil, in meq 100 g-1, for the horizon. The range of valid entries is from 0 to 400 meq 100 g-1, and tenths (one decimal place) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.100. 618.20 Electrical Conductivity Definition.—“Electrical conductivity” is the electrolytic conductivity of an extract from saturated soil paste. Classes.—The classes of salinity are: Salinity Class Electrical Conductivity (mmhos cm-1) Nonsaline 0 to <2 Very slightly saline 2 to <4 Slightly saline 4 to <8 Moderately saline 8 to <16 Strongly saline ≥16 Significance.—Electrical conductivity is a measure of the concentration of water-soluble salts in soils. It is used to indicate saline soils. High concentrations of neutral salts, such as sodium chloride and sodium sulfate, may interfere with the absorption of water by plants because the osmotic pressure in the soil solution is nearly as high as or higher than that in the plant cells. Salts may also interfere with the exchange capacity of nutrient ions, thereby resulting in nutritional deficiencies in plants. Electrical conductivity in the extract from a saturated paste is used for soil classification in the required characteristics for the salic horizon and in criteria for certain taxa such as Dystric great groups and Halic subgroups of Vertisols. Measurement.—The electrolytic conductivity of a saturated extract is the standard measure used to express salinity. Units of measure are decisiemens per meter (dS m-1), which are equivalent to millimhos per centimeter (mmhos cm-1), at 25 degrees C. The laboratory procedure used to measure electrical conductivity is described in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Estimates.—Field estimates of salinity are made from observations of visible salts on faces of peds, throughout the horizon matrix, on the soil surface, or some combination of the three; from plant growth or productivity; from the presence of native plant indicator species; and from field salinity meters. The occurrences of bare spots, salt-tolerant plants, and uneven crop growth are used as indicators of salinity and high electrical conductivity. When keyed to measurements, these observations help to estimate the amount of salts. Entries.—Enter the high, low, and representative values for the range of electrolytic conductivity of the saturation extract during the growing season for each horizon. If laboratory measurements or accurate field estimates are available, the high and low values do not need to correspond with salinity class limits. However, if data is limited, use the following ranges to represent the high and low values of the salinity classes: 0-2, 2-4, 4-8, 8-16, and 16-32 (or a reasonable high value for the strongly saline class) or use a combination of classes (for example, 2-8 for the high and low values). The range of valid entries is from 0 to 15000 mmhos cm-1, and tenths (one decimal place) are allowed. 618.21 Electrical Conductivity 1:5 (volume) Definition.—“Electrical conductivity 1:5 (volume)” is the electrolytic conductivity of of a diluted, unfiltered supernatant of 1 part soil to 5 parts distilled water as measured by volume. The measurement is only used for subaqueous soils. Classes.—See the salinity classes described above in the section “Electrical Conductivity.” The traditional salinity classes were designed as phase criteria for terrestrial soils and are not applicable to subaqueous soils. Significance.—Electrical conductivity (EC) 1:5 (volume) is a measure of the concentration of water-soluble salts in soils. It is used to indicate the threshold between freshwater and salt and brackish water subaqueous soils. Measuring EC in this manner is the best approach for subaqueous soils as samples containing reduced sulfide must be kept moist to avoid oxidation and production of sulfate slats that can increase the electrical conductivity. Salinity tolerance in plants is a measure of diminished plant growth at a threshold of 10% reduction in biomass. This is not the same as the maximum salinity tolerance which is an LD50 response. Electrical conductivity 1:5 (volume) is used for soil classification in the criteria for the Frasiwassents and Frasiwassists great groups. Measurement.—EC 1:5 (volume) measured in an unfiltered supernatant is the standard measure used to express salinity in subaqueous soils. EC 1:5 (volume) must be measured in a fresh, field wet sample (moisture content at sample collection) that has been refrigerated or even frozen because sulfides may oxidize during drying forming sulfate salts, which can increase the EC value. This method assumes that the salts in subaqueous soils are highly soluble chloride and sulfate salts, in a dissolved state, with no important contributions from minerals such as gypsum. Units of measure are decisiemens per meter (dS m-1), at 25 degrees C. The laboratory procedure used to measure electrical conductivity is described in an addendum to the Soil Survey Field and Laboratory Methods Manual Report No. 51, Soil Survey Field and Laboratory Methods Manual, Version 1.0, 2009, USDA, NRCS. Estimates.—Field estimates of salinity are made for subaqueous soils from observations of the presence of native plant indicator species and from measuring the water column with field salinity refractometers. Caution should be used in comparing water column salinity to soil salinity. Ground water discharge can decrease soil salinity and seasonal evaporation of seawater in barrier salt marshes can produce brine that sinks through the groundwater to collect in subsurface coarse-textured lenses. Salinty distributions in mainland associated soils tend to have a systematic decrease with depth while salinity in other subaqueous soils remain high with depth. Entries.—Enter the high, low, and representative values for the range of electrolytic conductivity 1:5 (volume) of the unfiltered supernatant for each horizon. The range of valid entries is from 0 to 100 dS m-1, and tenths (one decimal place) are allowed. 618.22 Elevation Definition.—“Elevation” is the vertical distance from mean sea level to a point on the earth’s surface. Significance.—Elevation, or local relief, influences the genesis of natural soil bodies. Elevation also may affect soil drainage within a landscape, salinity or sodicity within a climatic area, or soil temperature. Estimates.—Elevation is normally obtained from U.S. Geological Survey topographic maps or measured using altimeters or global positioning systems. Entries.—Enter the high, low, and representative values for elevation in meters for each map unit component. The range of valid entries is from -300 to 8550 meters, and tenths (one decimal place) are allowed. 618.23 Engineering Classification AASHTO Group Classification Definition.—“AASHTO group classification” is a system that classifies soils specifically for geotechnical engineering purposes that are related to highway and airfield construction. It is based on particle-size distribution and Atterberg limits, such as liquid limit and plasticity index. This classification system is covered in Standard No. M 145-82, published by the American Association of State Highway and Transportation Officials (AASHTO), and consists of a symbol and a group index. The classification is based on that portion of the soil that is smaller than 3 inches in diameter. Classes.—The AASHTO classification system identifies two general classifications: granular materials having 35 percent or less, by weight, particles smaller than 0.074 mm in diameter and silt-clay materials having more than 35 percent, by weight, particles smaller than 0.074 mm in diameter. These two divisions are further subdivided into seven main group classifications. Part 618, Subpart B, Exhibits, Section 618.83 shows the criteria for classifying soil in the AASHTO classification system. The group and subgroup classifications are based on estimated or measured grain-size distribution and on liquid limit and plasticity index values. Significance.—The group and subgroup classifications of this system aid in the evaluation of soils for highway and airfield construction. The classifications can help to make general interpretations relating to performance of the soil for engineering uses, such as highways and local roads and streets. Measurements.—Measurements involve sieve analyses for the determination of grain-size distribution of that portion of the soil between a 3 inch and 0.074 mm particle size. ASTM Designations D 422, C 136, and C 117 have applicable procedures for the determination of grain-size distribution. The liquid limit and plasticity index values (ASTM Designation D 4318) are determined for that portion of the soil having particles smaller than 0.425 mm in diameter (no. 40 sieve). Measurements, such as laboratory tests, are made on most benchmark soils and on other representative soils in survey areas. Estimates.—During soil survey investigations and field mapping activities, the soil is classified by field methods. This classification involves making estimates of particle-size fractions by a percentage of the total soil, minus the greater-than-3-inch fraction. Estimates of liquid limit and plasticity index are based on clay content and mineralogy relationships. Estimates are expressed in ranges that include the estimating accuracy as well as the range of values for the taxon. Definition.—The AASHTO group and subgroup classifications may be further modified by the addition of a group index value. The empirical group index formula was devised for approximate within-group evaluation of the “clayey granular” materials and the “silty-clay” materials. Significance.—The group index (GI) aids in the evaluation of the soils for highway and airfield construction. The index can help to make general interpretations relating to performance of the soil for engineering uses, such as highways and local roads and streets. In calculating the group index of A-2-6 and A-2-7 subgroups, only the PI portion of the formula is used. Entries.—The group index is reported to the nearest integer. If the calculated group index is negative, the group index value is zero. The minimum group index value is 0 and the maximum is 120. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.98. Unified Soil Classification Definition.—The “Unified soil classification system” is a system that classifies mineral and organic mineral soils for engineering purposes based on particle-size characteristics, liquid limit, and plasticity index. Classes The Unified soil classification system identifies three major soil divisions: Coarse grained soils having less than 50 percent, by weight, particles smaller than 0.074 mm in diameter. Highly organic soils that demonstrate certain organic characteristics. These divisions are further subdivided into a total of 15 basic soil groups. The major soil divisions and basic soil groups are determined on the basis of estimated or measured values for grain-size distribution and Atterberg limits. ASTM Designation D 2487 shows the criteria chart used for classifying soil in the Unified system, the 15 basic soil groups of the system, and the plasticity chart for the system. Significance.—The various groupings of this classification have been devised to correlate in a general way with the engineering behavior of soils. This correlation provides a useful first step in any field or laboratory investigation for engineering purposes. It can be used to make some general interpretations relating to probable performance of the soil for engineering uses. Measurement.—The methods for measurement are provided in ASTM Designation D 2487. Measurements involve sieve analysis for the determination of grain-size distribution of that portion of the soil between 3 inches and 0.074 mm in diameter (no. 200 sieve). ASTM Designations D 422, C 136, and C 117 have applicable procedures that are used, where appropriate, for the determination of grain-size distribution. Values for the Atterberg limits (liquid limit and plasticity index) are also used. Specific tests are made for that portion of the soil having particles smaller than 0.425 mm in diameter (no. 40 sieve) according to ASTM Designations D 423 and D 424. Measurements, such as laboratory tests, are made on most benchmark soils and on other representative soils in survey areas. Estimates.—The methods for estimating are provided in ASTM Designation D 2488. During all soil survey investigations and field mapping activities, the soil is classified by field methods. The methods include making estimates of particle-size fractions by a percentage of the total soil. The Atterberg limits are also estimated based on the wet consistency, ribbon or thread toughness, and other simple field tests. These tests and procedures are explained in ASTM Designation D 2488. If samples are later tested in the laboratory, adjustments are made to field procedures as needed. Estimates are expressed in ranges that include the estimating accuracy as well as the range of values from one location to another within the map unit. If an identification is based on visual-manual procedures, it must be clearly stated so in reporting. 618.24 Erosion Accelerated, Kind Definition.—“Erosion accelerated, kind” is the type of detachment and removal of surface soil particles that is largely affected by human activity. Significance.—The type of accelerated erosion is important in assessing the current health of the soil and in assessing its potential for different uses. Erosion, whether natural or induced by humans, is an important process that affects soil formation and may remove all or parts of the soils formed in the natural landscape. Classes.—There are five kinds of accelerated erosion: Water erosion, sheet Water erosion, rill Water erosion, gully Water erosion, tunnel Wind erosion Entries.—Enter the appropriate class for each map unit component. Multiple entries are allowable, but a representative value should be indicated. 618.25 Erosion Class Definition.—“Erosion class” is the class of accelerated erosion. Significance The degree of erosion that has taken place is important in assessing the health of the soil and in assessing the soil’s potential for different uses. Erosion is an important process that affects soil formation and may remove all or parts of the soils formed in natural landscapes. Removal of increasing amounts of soil increasingly alters various properties and capabilities of the soil. Properties and qualities affected include bulk density, organic matter content, tilth, and water infiltration. Altering these properties affects the productivity of the soil. Estimates.—During soil examinations, estimate the degree to which soils have been altered by accelerated erosion. The Soil Survey Manual describes the procedures involved. Classes.—There are five erosion classes: None - deposition Class 1 Class 2 Class 3 Class 4 Entries.—Enter the appropriate class for each map unit component. 618.26 Excavation Difficulty Classes Definition.—“Excavation difficulty classes” are used for soil layers, horizons, pedons, or geologic layers and estimate the difficulty of making an excavation into them. Excavation difficulty, in most instances, is strongly controlled by water state, which should be specified. Classes.—The excavation difficulty classes are defined in the following table. Class Definition Low Excavations can be made with a spade using arm-applied pressure only. Neither application of impact energy nor application of pressure with the foot to a spade is necessary. Moderate Arm-applied pressure to a spade is insufficient. Excavation can be accomplished quite easily by application of impact energy with a spade or by foot pressure on a spade. High Excavation with a spade can be accomplished with difficulty. Excavation is easily possible with a full length pick, using an over-the-head swing. Very High Excavation with a full length pick, using an over-the-head swing, is moderately to markedly difficult. Excavation is possible in a reasonable period of time with a backhoe mounted on a 40 to 60 kW (50-80 hp) tractor. Extremely High Excavation is nearly impossible with a full length pick using an over-the-head arm swing. Excavation cannot be accomplished in a reasonable time period with a backhoe mounted on a 40 to 60 kW (50-80 hp) tractor. Significance.—Excavation difficulty classes are important for evaluating the cost and time needed to prepare shallow excavations. Estimates.—Estimates of excavation difficulty classes are made from field observations. Entries.—Enter the appropriate class for each horizon. The allowable entries are low, moderate, high, very high, and extremely high. 618.27 Exchangeable Sodium Definition.—“Exchangeable sodium” is a measure of soil exchangeable sodium ions that may become active by cation exchange. It is the fraction of the cation-exchange capacity of a soil that is occupied by sodium ions, expressed as a percentage. Significance.—Exchangable sodium percentage (ESP) is used for soil classification in the required characteristics for the natric horizon, in the key to soil orders and key to suborders of Inceptisols and Mollisols, and in criteria for certain taxa such as Sodic subgroups. Soils that have values for exchangeable sodium of 15 percent or more may have an increased dispersion of organic matter and clay particles, reduced saturated hydraulic conductivity and aeration, and a general degradation of soil structure. Measurement.—The ESP is calculated by several methods which use the results of separate procedures to measure the sodium extractable by NH4OAc and the cation-exchange capacity by NH4OAc, pH 7.0 (CEC-7) as outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Units of measure for extractable sodium and cation-exchange capacity are centimoles per kilogram (cmol(+) kg-1) of soil, which are equivalent to milliequivalents per 100 grams (meq 100 g-1) of soil. In some soils with high salt content (i.e., >20 dS m-1) the ESP is calculated using the sodium extractable by NH4OAc, the cation-exchange capacity by NH4OAc, pH 7.0 (CEC-7), the water saturation percentage, and the water-soluble sodium (mmol (+) L-1). Entries.—Enter high, low, and representative values as percentages for each horizon for which data is available. The range of valid entries is from 0 to 100 percent, and only whole numbers (integers) are allowed. 618.28 Extractable Acidity Definition.—“Extractable acidity” is a measure of soil exchangeable hydrogen ions that may become active by cation exchange. Significance.—Extractable acidity is important for certain evaluations of soil nutrient availability or of the effect of waste additions to the soil. Extractable acidity is indirectly important data for soil classification because it is needed to calculate cation-exchange capacity by sum of cations (at pH 8.2). The cation-exchange capacity by the sum of cations method is used to calculate percent base saturation by sum of cations. Measurement.—Extractable acidity is determined by a method using a solution of barium chloride-triethanolamine buffered at pH 8.2, as outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Units of measure are centimoles per kilogram (cmol(+) kg-1) of soil, which are equivalent to milliequivalents per 100 grams (meq 100 g-1) of soil. Entries.—Enter the range of extractable acidity in milliequivalents per 100 grams (meq 100 g-1) of soil for the horizon. The range of valid entries is from 0 to 250 meq 100 g-1, and tenths (one decimal place) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.101. 618.29 Extractable Aluminum Definition.—“Extractable aluminum” is the amount of aluminum that approximates the aluminum considered exchangeable. It is a measure of the “active” acidity present in soils with a 1:1 water pH ≤5.5. Significance.—Extractable aluminum is important for certain evaluations of soil nutrient availability and of toxicities. An aluminum saturation of about 60 percent is usually regarded as toxic to most plants. It may be a useful measurement for assessing potential lime needs for acid soils. Extractable aluminum is used for soil classification in the criteria for Alic and some Eutric subgroups of Andisols. Measurement.—Extractable aluminum is determined by a method using a solution of one normal potassium chloride, as outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Units of measure are centimoles per kilogram (cmol(+) kg-1) of soil, which are equivalent to milliequivalents per 100 grams (meq 100 g-1) of soil. Entries.—Enter the range of extractable aluminum as milliequivalents per 100 grams (meq 100 g-1) of soil for the horizon. The range of valid entries is from 0 to 150 meq 100 g-1, and hundredths (two decimal places) are allowed. 618.30 Flooding Frequency Class, Duration Class, and Month Definition.—“Flooding” is the temporary covering of the soil surface by flowing water from any source, such as streams overflowing their banks, runoff from adjacent or surrounding slopes, inflow from high tides, or any combination of sources. Shallow water standing or flowing that is not concentrated as local runoff during or shortly after rain or snowmelt is excluded from the definition of flooding. Chapter 3 of the Soil Survey Manual provides additional information. Standing water (ponding) or water that forms a permanent covering is also excluded from the definition. Classes.—Estimates of flooding class are based on the interpretation of soil properties and other evidence gathered during soil survey fieldwork. Flooding hazard is expressed by flooding frequency class, flooding duration class, and time of year that flooding occurs. Not considered here, but nevertheless important, are velocity and depth of floodwater. Frequencies used to define classes are generally estimated from evidence related to the soil and vegetation. They are expressed in wide ranges that do not indicate a high degree of accuracy. Flooding frequencies that are more precise can be calculated by performing complex analyses used by engineers. The class “very frequent” is used for areas subject to daily and monthly high tides. Flooding Frequency Class.—Flooding frequency class indicates the number of times flooding occurs over a period of time. The classes of flooding are defined as follows: Flooding Frequency Class Definition None No reasonable possibility of flooding; one chance out of 500 of flooding in any year or less than 1 time in 500 years. Very rare Flooding is very unlikely but is possible under extremely unusual weather conditions; less than 1 percent chance of flooding in any year or less than 1 time in 100 years but more than 1 time in 500 years. Rare Flooding is unlikely but is possible under unusual weather conditions; 1 to 5 percent chance of flooding in any year or nearly 1 to 5 times in 100 years Occasional Flooding is expected infrequently under usual weather conditions; 5 to 50 percent chance of flooding in any year or 5 to 50 times in 100 years. Frequent Flooding is likely to occur often under usual weather conditions; more than a 50 percent chance of flooding in any year (i.e., 50 times in 100 years), but less than a 50 percent chance of flooding in all months in any year. Very frequent Flooding is likely to occur very often under usual weather conditions; more than a 50 percent chance of flooding in all months of any year. Flooding Duration Class.—The average duration of inundation per flood occurrence is given only for the occasional, frequent, and very frequent classes (defined above). Flooding Duration Class Duration Extremely brief 0.1 to 4 hours Very brief 4 hours to < 2 days Brief 2 days to < 7 days Long 7 days to < 30 days Very long > 30 days Assignment.—Yearly flooding frequency classes are assigned to months and indicate the months of occurrence and not the frequency of the flooding during the month, except for the very frequent class. The time period expressed includes two-thirds to three-fourths of the occurrences. Time period and duration of flooding are the most critical factors that determine the growth and survival of a given plant species. Flooding during the dormant season has few if any harmful effects on plant growth or mortality and may improve the growth of some species. If inundation from floodwater occurs for long periods during the growing season, the soil becomes oxygen deficient and plants may be damaged or killed. Significance.—The susceptibility of soils to flooding is an important consideration for building sites, sanitary facilities, and other uses. Floods may be less costly per unit area of farmland as compared to that of urban land, but the loss of crops and livestock can be disastrous. Estimates.—The most precise evaluation of flood-prone areas for stream systems is based on hydrologic studies. The area subject to inundation during a flood of a given frequency, such as one with a 1 percent or 2 percent chance of occurrence, is generally determined by one of two basic methods. The first method is used if stream flow data are available. In this method, the data are analyzed to determine the magnitude of floods of different frequencies. Engineering studies are made to determine existing channel capacities and flow on the flood plain by the use of valley cross sections and water-surface profiles. The second method is used if stream flow data are not available. In this method, hydrologists make an estimate of flood potential from recorded data on rainfall. They consider such factors as— Size, slope, and shape of the contributing watershed. Hydrologic characteristics of the soil. Land use and treatment. Hydraulic characteristics of the valley and channel system. With the use of either method, soil surveys can aid in the delineation of flood-prone areas. Possible sources of flooding information are— NRCS project-type studies, such as those arising from Public Law 556, flood protection, river basin, or resource conservation and development projects. Flood hazard analyses. Corps of Engineers (COE) flood plain information reports. Special flood reports. Local flood-protection and flood-control project reports. Department of Housing and Urban Development flood-insurance study reports. Maps by the U.S. Geological Survey (USGS), NRCS, Tennessee Valley Authority, COE, or National Oceanic and Atmospheric Administration. General estimates of flooding frequency and duration are made for each soil. However, in intensively used areas where construction has materially altered the natural water flow, flood studies are needed to adequately reflect present flooding characteristics. Soil scientists collect and record evidence of flood events during the course of the soil survey. The extent of flooded areas, flood debris in trees, damage to fences and bridges, and other signs of maximum water height are recorded. Information that is helpful in delineating soils that have a flood hazard is also obtained. Hydrologists may have flood stage predictions that can be related to kinds of soil or landscape features. Conservationists and engineers may have recorded elevations of high flood marks. Local residents may have recollections of floods that can help to relate the events to kinds of soil, topography, and geomorphology. Certain landscape features have developed as the result of past and present flooding and include former river channels, oxbows, point bars, alluvial fans, meander scrolls, sloughs, natural levees, backswamps, sand splays, and terraces. Most of these features are easily recognizable on aerial photographs when the photo image is compared to on-the-ground observations. Different kinds of vegetation and soils are normally associated with these geomorphic features. The vegetation that grows in flood areas may furnish clues to past flooding. In the central and southeastern United States, the survival of trees in flood-prone areas depends on the frequency, duration, depth, and time of flooding and on the age of the tree. Past flooding may sometimes leave clues in the soil, such as— Thin strata of material of contrasting color, texture, or both. An irregular decrease in organic matter content, not due to human-alteration by mixing or transportation of material, which is an indication of a buried genetic surface horizon. Soil layers that have abrupt boundaries to contrasting kinds of material, which indicate that the materials were laid down suddenly at different times and were from different sources or were deposited from stream flows of different velocities. Artifacts which are easily moved and deposited by flood waters (e.g., plastic bottles). Entries.—If a map unit component floods, then the annual flooding frequency and duration are populated, as stated above in Section 618.30 B, for the specific months in which the flooding events most commonly occur. All other months have records in the Component Month table but the data elements for frequency and duration are left NULL. Flooding entries reflect the current existing and mapped condition with consideration for dams, levees, and other human-induced changes affecting flooding frequency and duration. Enter the flooding duration class name that most nearly represents the soil: extremely brief, very brief, brief, long, or very long. 618.31 Fragments in the Soil Definition.—“Fragments” are unattached, cemented pieces of bedrock, bedrock-like material, durinodes, concretions, nodules, or pedogenic horizons (e.g., petrocalcic fragments) 2 mm or larger in diameter and unprocessed woody material 20 mm or larger in diameter in organic soils. Fragments are separated into three types: rock fragments, pararock fragments (which are distinguished by cementation class), and wood fragments. The words “rock” and “pararock” are used here in the broad sense and do not connote only natural fragments of geologic material. Some artifacts behave in a similar manner to fragments in the soil. See Section 618.5 for detailed information on the measurement, classes, and data entries for artifacts. Rock fragments are unattached pieces of geologic or pedogenic material 2 mm in diameter or larger that are strongly cemented or more resistant to rupture. Rock fragments from 2 mm to 75 mm (3 inches) are considered when estimating the percent passing sieves, as discussed in section 618.47. Pararock fragments are unattached, cemented pieces of geologic or pedogenic material 2 mm in diameter or larger that are extremely weakly cemented to moderately cemented. These fragments are not retained on sieves because they are crushed by grinding during sample preparation. Wood fragments are unprocessed (i.e., naturally occurring) woody materials that cannot be crushed between the fingers when moist or wet and are larger than 20 mm in size. Wood fragments are only used in organic soils. They are comparable to rock and pararock fragments in mineral soils. Processed wood products, whether treated or untreated, are considered artifacts and not wood fragments. Significance The fraction of the soil 2 mm or larger in diameter has an impact on the behavior of the whole soil. Soil properties, such as available water capacity, cation-exchange capacity, saturated hydraulic conductivity, structure, and porosity, are affected by the volume, composition, and size distribution of rock fragments in the soil. Fragments also affect the management of the soil and are used as interpretation and classification criteria (e.g., particle-size and substitute classes). Terms related to volume, size, and hardness of fragments are used as texture modifier terms. Generally, the fraction of soil greater than 75 mm (3 inches) in diameter is not included in the engineering classification systems. However, it can be added as a descriptive term to the group name, for example, poorly graded gravel with silt, sand, cobbles, and boulders. Estimates of the percent of cobbles and boulders are presented in the soil descriptions for a group name (see ASTM Designation D 2487, appendix X1.1). A small amount of these larger particles generally has little effect on soil properties. The particles may, however, have an effect on the use of a soil in certain types of construction. Often, the larger portions of a soil must be removed before the material can be spread in thin layers, graded, or compacted and graded to a smooth surface. As the quantity of this “oversized” fraction increases, the properties of the soil can be affected. If the larger particles are in contact with each other, the strength of the soil is very high and the compressibility very low. If voids exist between the larger particles, the soil will likely have high saturated hydraulic conductivity and may undergo some internal erosion as a result of the movement of water through the voids. Most of the smaller and more rapid construction equipment normally used in excavating and earthmoving cannot be used if the oversize fraction of a soil is significant. Measurement.— The fraction from 2 to 75 mm in diameter may be measured in the field. However, 50 to 60 kg of sample material may be necessary if there is an appreciable amount of fragments near 75 mm. An alternative means of measuring is to visually estimate the volume of the 20- to 75-mm fraction, then sieve and weigh the 2- to 20-mm fraction. The fraction 75 mm (3 inches) or greater is usually not included in soil samples taken in the field for laboratory testing. Measurements can be made in the field by weighing the dry sample and the portion retained on a 3-inch screen. The quantity is expressed as a weight percentage of the total soil. A sample as large as 200 pounds to more than a ton may be needed to assure that the results are representative. Measurements of the fraction from 75 to 250 mm (3 to 10 inches) and the fraction greater than 250 mm (10 inches) in diameter are usually obtained from volume estimates. Estimates Estimates are usually made by visual means and are on the basis of percent by volume. The percent by volume is converted to percent by weight by using the average bulk unit weights for soil and rock. These estimates are made during investigation and mapping activities in the field. They are expressed as ranges that include the estimating accuracy as well as the range of values for a component. Measurements or estimates of fragments less than strongly cemented are made prior to any rolling or crushing of the sample. Rock Fragments Greater Than 10 Inches (250 mm) Definition.—“Rock fragments greater than 10 inches” is the percent by weight of the horizon occupied by rock fragments greater than 10 inches (250 mm) in diameter. Although the upper limit is undefined, for practical purposes it generally is no larger than a pedon, up to 10 meters square. For flat rock fragments, the intermediate dimension is used for the 250 mm (10 inch) measurement. For example, a flat-shaped rock fragment that is 100 mm x 250 mm x 380 mm has an intermediate dimension of 250 mm and is not counted as greater than 250 mm. A flat-shaped rock fragment that is 100 mm x 275 mm x 380 mm has an intermediate dimension of 275 mm and is counted as greater than 250 mm. Entries.—Enter the high, low, and representative values in the Component Horizon Table in the NASIS database as whole number percentages for each horizon, as appropriate. Rock Fragments 3 to 10 Inches (75 to 250 mm) Definition.—“Rock fragments 3 to 10 inches” is the percent by weight of the horizon occupied by rock fragments 3 to 10 inches (75 to 250 mm) in diameter. Entries.—Enter the high, low, and representative values in the Component Horizon Table in the NASIS database as whole number percentages for each horizon, as appropriate. Fragment Kind Definition.—“Fragment kind” is the lithology or composition of the 2 mm or larger fraction of the soil (20 mm or larger for wood fragments). Entries.—Enter the appropriate fragment kind name for the record of fragments populated in the Component Horizon Fragments Table in the NASIS database. The class names appear in a choice list and can also be viewed in the NASIS data dictionary. Fragment Roundness Definition.—“Fragment roundness” is an expression of the sharpness of edges and corners of fragments. Entries.—Enter the appropriate fragment roundness class name for the record of fragments populated in the Component Horizon Fragments Table in the NASIS database. Fragment Hardness Definition.—“Fragment hardness” is equivalent to the “rupture resistance cemented” of a fragment of specified size that has been air dried and then submerged in water. Measurements.—Measurements are made using the procedures and classes of cementation that are listed with the rupture resistance classes in the Soil Survey Manual. Classes are described for block-like specimens about 25-30 mm on edge, which are air dried and then submerged in water for at least 1 hour. The specimen is compressed between an extended thumb and forefinger, between both hands, or between a foot and a nonresilient flat surface. If the specimen resists compression, a weight is dropped onto it from progressively greater heights until it ruptures. Failure is considered at the initial detection of deformation or rupture. Stress applied in the hand should be over a 1-second period. The tactile sense of the class limits may be learned by applying force to top-loading scales and sensing the pressure through the tips of the fingers or through the ball of the foot. Postal scales may be used for the resistance range that is testable with the fingers. A bathroom scale may be used for the higher rupture resistance range. Significance.—The hardness of a fragment is significant where the rupture resistance class is strongly cemented or greater. These classes can impede or restrict the movement of soil water vertically through the soil profile and have a direct impact on the quality and quantity of ground water and surface water. Classes.—The fragment hardness (rupture resistance) classes are— Extremely weakly cemented Very weakly cemented Weakly cemented Strongly cemented Very strongly cemented Indurated Entries.—Enter the appropriate class namefor each record of fragments populated in the Component Horizon Fragments Table in the NASIS database. Choose the term without the word “cemented” (i.e., choose “moderately” to represent the moderately cemented class). Fragment Shape Definition.—“Fragment shape” is a description of the overall shape of the fragment. Significance.—Fragment shape is important for fragments that are too large to be called channers or flagstones. Classes.—The fragment shape classes are flat and nonflat. Entries.—Enter the appropriate fragment shape class name for each record of fragments populated in the Component Horizon Fragments Table in the NASIS database. Fragment Size Definition.—“Fragment size” is based on the multiaxial dimensions of the fragment. Significance.—The size of fragments is significant to the use and management of the soil. Fragment size is used as a criterion in naming map units. It affects equipment use, excavation, construction, and recreational uses. Classes.—Classes of fragment size are subdivided as flat or nonflat based on the shape of the fragments (described above). Flat fragment classes are: Flat Fragment Class Length of Fragment (mm) Channers 2-150 Flagstones 150-380 Stones 380-600 Boulders ≥600 Nonflat fragment classes are: NonFlat Fragment Class Diameter (mm) Gravel 2-75 fine gravel 2-5 medium gravel 5-20 coarse gravel 20-75 Cobbles 75-250 Stones 250-600 Boulders ≥600 Gravel is a collection of fragments having a diameter ranging from 2 to 75 mm. Individual fragments in this size range are properly referred to as pebbles, not “gravels.” For fragments that are less than strongly cemented, “para” is used as a prefix to the above terms (e.g., paracobbles). Entries.—Enter the high, low, and representative values of each size class populated in the Component Horizon Fragments Table in the NASIS database. The range of valid entries is from 2 to 3,000 millimeters, and only whole numbers (integers) are allowed. Fragment Volume Definition.—“Fragment volume” is the volume percentage of the horizon occupied by the 2 mm or larger fraction (20 mm or larger for wood fragments) on a whole soil base. Significance.—The volume occupied by the 2 mm or larger fraction is important in selecting texture modifiers (i.e., gravelly, very gravelly, extremely paragravelly). Entries.—Enter the high, low, and representative values for the percent volume present of each size class and kind of fragment populated in the Component Horizon Fragments Table in the NASIS database. The range of valid entries is from 0 to 100 percent, and only whole numbers (integers) are allowed. 618.32 Free Iron Oxides Definition.—“Free iron oxides” are secondary iron oxides, such as goethite, hematite, ferrihydrite, lepidocrocite, and maghemite. These forms of iron may occur as discrete particles, as coatings on other soil particles, or as cementing agents between soil mineral grains. They consist of iron extracted by dithionite-citrate from the fine-earth fraction. Significance.—The amount of iron that is extractable by dithionite-citrate is used in the ferritic, ferruginous, parasesquic, and sesquic mineralogy classes defined in Soil Taxonomy. The ratio of dithionite-citrate (free) iron to total iron in a soil is a measure of the degree of soil weathering. Free iron oxides are important in the soil processes of podzolization and laterization and play a significant role in the phosphorous fixation ability of soils. Measurement.— Free iron oxides are measured as the amount extracted by a solution of sodium dithionite and sodium-citrate using a method outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Entries.—Enter high, low, and representative values as percentages for each horizon for which data is available. The range of valid entries is from 0 to 100 percent, and hundredths (two decimal places) are allowed. 618.33 Frost Action, Potential Definition.—“Potential frost action” is a rating of the susceptibility of the soil to upward or lateral movement by the formation of segregated ice lenses. It rates the potential for frost heave and the subsequent loss of soil strength when the ground thaws. Soils are susceptible to the formation of ice lenses, which results in frost heave and subsequent loss of soil strength. High Soils are highly susceptible to the formation of ice lenses, which results in frost heave and subsequent loss of soil strength. Significance.—Damage from frost action results from the formation of segregated ice crystals and ice lenses in the soil and the subsequent loss of soil strength when the ground thaws. Frost heave damages highway and airfield pavements. It is less of a problem for dwellings and buildings that have footings which extend below the depth of frost penetration. In cold climates, unheated structures that have concrete or asphalt floors can be damaged by frost heave. Driveways, patios, and sidewalks can heave and crack. The thawing of the ice causes a collapse of surface elevation and produces free-water perches on the still-frozen soil below. Soil strength is reduced. Backslopes and side slopes of cuts and fills can slough during thawing. Seedlings and young plants of clover, alfalfa, wheat, and oats can be raised out of the soil or have their root systems damaged by frost heave. Estimates Freezing temperatures, soil moisture, and susceptible soils are needed for the formation of segregated ice lenses. Ice crystals begin to form in the large pores first. Water in small pores or water that was adsorbed on soil particles freezes at lower temperatures. This supercooled water is strongly attracted to the ice crystals, moves toward them, and freezes on contact with them. The resulting ice lens continues to grow in width and thickness until all available water that can be transported by capillary has been added to the ice lens and a further supply cannot be made available because of the energy requirements. Soil temperatures must drop below 0° C for frost action to occur. Generally, the more slowly and deeply the frost penetrates, the thicker the ice lenses are and the greater the resulting frost heave is. Part 618, Subpart B, Exhibits, Section 618.85 is a map that shows the design freezing index values in the continental United States. The values are the number of degree days below 0° C for the coldest year in a period of 10 years. The values indicate duration and intensity of freezing temperatures. The 250 isoline is the approximate boundary below which frost action ceases to be a problem. Except on the West Coast, the frost action boundary corresponds closely to the functional boundary between the mesic and thermic soil temperature regimes as defined in Soil Taxonomy. More information is provided in the U.S. Army Engineer School, Student Reference, 1967, Soil Engineering, Section I, Volume II, Chapters VI-IX, Fort Belvoir, Virginia. Water necessary for the formation of ice lenses may come from a high water table or from infiltration at the surface. Capillary water in voids and adsorbed water on particles also contribute to ice lens formation but, unless this water is connected to a source of free water, the amount generally is insufficient to produce significant ice segregation and frost heave. The potential intensity of ice segregation is dependent to a large degree on the effective soil pore size and soil saturated hydraulic conductivity, which are related to soil texture. Ice lenses form in soils in which the pores are fine enough to hold quantities of water under tension but coarse enough to transmit water to the freezing front. Soils that have a high content of silt and very fine sand have this capacity to the greatest degree and hence have the highest potential for ice segregation. Clayey soils hold large quantities of water but have such slow saturated hydraulic conductivity that segregated ice lenses are not formed unless the freezing front is slow moving. However, sandy soils have large pores and hold less water under lower tension. As a result, freezing is more rapid and the large pores permit ice masses to grow from pore to pore, entombing the soil particles. Thus, in coarse grained soils, segregated ice lenses are not formed and less displacement can be expected. Estimates of potential frost action generally are made for soils in mesic or colder temperature regimes. Exceptions are on the West Coast, where the mesic-thermic temperature line crosses below the 250 isoline, as displayed in Part 618, Subpart B, Exhibits, Section 618.85, and along the East Coast, where the soil climate is moderated by the ocean. Mesic soils that have a design freezing index of less than 250 degree days should not be rated because frost action is not likely to occur. The estimates are based on bare soil that is not covered by insulating vegetation or snow. They are also based on the moisture regime of the natural soil. The ratings can be related to manmade modifications of drainage or to irrigation systems on an onsite basis. Frost action estimates are made for the whole soil to the depth of frost penetration, to bedrock, or to a depth of 2 meters (6.6 feet), whichever is shallowest. Part 618, Subpart B, Exhibits, Section 618.84 is a guide for making potential frost action estimates. It uses the soil moisture regimes and taxonomic family particle-size classes as defined in Soil Taxonomy. Entries.—Enter one of the following classes for the whole soil: low, moderate, or high. If frost action is not a problem, enter “none.” 618.34 Gypsum Definition.—“Gypsum” is the percent, by weight, of hydrated calcium sulfates in the <20 mm fraction of soil. Significance.—Gypsum is partially soluble in water and can be dissolved and removed by water. Soils with more than 10 percent gypsum may collapse if the gypsum is removed by percolating water. Gypsum is corrosive to concrete. Corrosion of concrete is most likely to occur in soils that are more than about 1 percent gypsum when wetting and drying occurs. Gypsum percentage is used for soil classification in the required characteristics for gypsic and petrogypsic horizons, the gypseous substitute classes, several strongly contrasting particle-size classes, and the hypergypsic, gypsic, and carbonatic mineralogy classes. Entries.—Enter the high, low, and representative values to represent the range in gypsum content as a weight percent of the soil fraction less than 20 mm in size. Round values to the nearest 5 percent for layers that are more than 5 percent gypsum and to the nearest 1 percent for layers that are less than 5 percent gypsum (for example, 0-1, 1-5, 5-10). If the horizon does not have gypsum, enter “0.” The range of valid entries is from 0 to 120 percent, and only whole numbers (integers) are allowed. 618.35 Horizon Depth to Bottom Definition.—“Horizon depth to bottom” is the distance from the top of the soil to the base of the soil horizon. Measurement.— Distance is measured from the top of the soil, which is defined as the top of the mineral soil, or, for soil with “O” horizons, the top of any organic layer that is at least slightly decomposed. For soils that are covered by 80 percent or more rock or pararock fragments, the top of the soil is the mean height of the top of the fragments. See chapter 3 in the Soil Survey Manual for a complete discussion. Measurement should be estimated to a depth of 200 cm for most soils and to a depth at least 25 cm below a lithic contact if the contact is above 175 cm. For soils, including those that have a root-restricting contact such as a paralithic contact, the lowest horizon bottom should extend to a depth of at least 25 cm below the contact or to a depth of 200 cm, whichever is shallower. Entries.—Horizon depth to bottom values used to populate component data in NASIS are not specific to any one point but rather are a reflection of commonly observed values based on field observations and are intended to model the component as it occurs throughout the map unit. Enter the high, low, and representative values in whole centimeters. The high value represents either the greatest depth to which the base of the horizon extends or, for horizons extending beyond the limit of field observation, is the depth to which observation was made (usually no more than 200 cm but at least 150 cm). 618.36 Horizon Depth to Top Definition.—“Horizon depth to top” is the distance from the top of the soil to the upper boundary of the soil horizon. Measurement.— Distance is measured from the top of the soil, which is defined as the top of the mineral soil, or, for soils with “O” horizons, the top of any organic layer that is at least slightly decomposed. For soils that are covered by 80 percent or more rock or pararock fragments, the top of the soil is the mean height of the top of the fragments. See chapter 3 in the Soil Survey Manual for a complete discussion. Entries.—Horizon depth to top values used to populate component data in NASIS are not specific to any one point but rather are a reflection of commonly observed values based on field observations and are intended to model the component as it occurs throughout the map unit. Enter the high, low, and representative values in whole centimeters. See Section 618.37, “Horizon Designation,” for a discussion as to how to list E/B and E and Bt type horizons. 618.37 Horizon Designation Definition.—“Horizon designation” is a concatenation of three kinds of symbols used in various combinations to identify layers of soil that reflect the investigator’s interpretations of genetic relationships among layers within a soil. Significance.—Soils vary widely in the degree to which horizons are expressed. The range is from little or no expression to strong expression. Layers of different kinds are identified by symbols. Designations are provided for layers that have been changed by soil formation and for those that have not. Designations are assigned after comparison of the observed properties of the layer with properties inferred for the material before it was affected by soil formation. Designations of genetic horizons express a qualitative judgment about the kind of changes that are believed to have taken place. A more detailed discussion can be reviewed in the latest edition of the Keys to Soil Taxonomy and in the Soil Survey Manual, Chapter 3. Horizon designations shown in field pedon descriptions (point data) represent a specific location on the landscape. Horizon designations used to populate component data in NASIS are not specific to any one point but rather are a reflection of commonly observed horizon sequences based on field observations and are intended to model the component as it occurs throughout the map unit so that accurate interpretations can be derived. Entries.—Enter combinations of symbols to reflect master horizons and their vertical subdivisions. Commonly occurring master horizon sequences, identified in field pedon descriptions (e.g., A-Bt1-Bt2-Btk-2Bk-2C), are used for soil components in NASIS. Generalized horizon layer designations (e.g., H1-H2-H3 etc.) may be used instead of genetic horizon designations. Users of generalized horizon layer designations must be cognizant of the fact that certain master horizon designations such as O, Cr, and R, are still required entries in NASIS for proper soil interpretations. It is not possible to include all master horizon sequences observed in individual pedon descriptions when populating the interpretive horizons of soil components. For example, if an Ap horizon is present and the former E horizon is incorporated into the Ap horizon, a sequence of Ap-E should not be used because the horizons would not normally occur together. Judgment is required when determining how much detail is required to represent the component adequately for interpretations. Care must be taken to maintain important differences between generic horizons whose range in properties are specified separately on the official series descriptions. Some general guidance follows. Master horizons (i.e., O, A, E, B) represent unique sets of pedogenic processes and therefore should not be combined when aggregating data, even if their basic properties are similar. Transitional horizons, such as EB and BC, should be recorded if they are commonly more than about 10 cm thick. After applying careful judgment, some horizons may be combined in order to avoid overly complex horizon sequences representing map unit components. Keep in mind, however, that once combined, any useful information gained by their separation is lost. Transitional horizons thinner than about 10 cm may be combined with an adjacent master horizon if they are not considered important to the interpretation of the component. Master horizon subdivisions showing genetic variations not deemed significant to interpretations may be combined. For example, a horizon sequence of Bt1-Bt2-Bt3, based solely on color variation and having no other significant differences, may be combined and shown simply as Bt. Master horizon subdivisions showing genetic variations that are deemed significant to interpretations should not be combined. For example, a horizon sequence of Bt-Btx-BC should not be combined. Do not combine horizons that straddle the criteria break to a diagnostic horizon. For example, a Bk1 with insufficient calcium carbonate content to qualify as a calcic horizon should not be combined with a Bk2 that is a calcic horizon. Enter only what the documentation can support. Combination horizons (E and Bt, Btn/E, E/Bt, etc.) should be entered as two separate horizon records, such as one for the E part of the horizon and the second for the Bt part of the horizon. Both records must have the same horizon designations assigned (e.g., “E/Bt”). But these separate horizon records must have different RV depth values for the top and bottom depths. The RV horizon depths must be completely in sync with no duplication, overlaps, or gaps. For example, the E part of a E/Bt horizon could have RV depths of 20 to 35 cm and the Bt part of the E/Bt horizon could have RV depths of 35 to 50 cm. The depth values for the “Low” and “High” columns of the horizon top and bottom depths may be populated to identify the overlapping nature of the horizon (e.g., both records may have the same low value for the top depth of 10 cm). Soil property data elements would be populated for each part to describe the characteristics of that separate part of the combination horizon. Allowable codes are listed in the NASIS data dictionary (see Part 618, Subpart B, Exhibits, Section 618.82 for the Web address of the current version). The rules for use of horizon designations are in the latest edition of the Keys to Soil Taxonomy and in the Soil Survey Manual,Chapter 3. 618.38 Horizon Thickness Definition.—“Horizon thickness” is a measurement from the top to bottom of a soil horizon throughout its areal extent. Measurement.— Soil horizon thickness varies on a cyclical basis. Measurements should be made to record the range in thickness as it normally occurs in the soil. Entries.—Horizon thickness values used to populate component data in NASIS are not specific to any one point but rather are a reflection of commonly observed values based on field observations and are intended to model the component as it occurs throughout the map unit. Enter the high, low, and representative values in whole centimeters. The minimum allowable entry is 1 cm. For horizons extending beyond the limit of field observation, thickness is only populated to the depth at which an observation was made. 618.39 Hydrologic Group Definition The complete definition and official criteria for hydrologic soil groups are available online in the National Engineering Handbook. “Hydrologic group” is a group of soils having similar runoff potential under similar storm and cover conditions. Soil properties that influence runoff potential are those that influence the minimum rate of infiltration for a bare soil after prolonged wetting and when not frozen. These properties are depth to a seasonal high water table, saturated hydraulic conductivity after prolonged wetting, and depth to a layer with a very slow water transmission rate. Changes in soil properties caused by land management or climate changes also cause the hydrologic soil group to change. The influence of ground cover is treated independently. Classes.—The soils in the United States are placed into four groups, A, B, C, and D, and three dual classes, A/D, B/D, and C/D. Significance.—Hydrologic groups are used in equations that estimate runoff from rainfall. These estimates are needed for solving hydrologic problems that arise in planning watershed-protection and flood-prevention projects and for planning or designing structures for the use, control, and disposal of water. Measurements.—The original classifications assigned to soils were based on the use of rainfall-runoff data from small watersheds and infiltrometer plots. From these data, relationships between soil properties and hydrologic groups were established. Estimates.— Assignment of soils to hydrologic groups is based on the relationship between soil properties and hydrologic groups. Wetness characteristics, water transmission after prolonged wetting, and depth to very slowly permeable layers are properties used in estimating hydrologic groups. 618.40 Landform Definition.—“Landform” is any physical, recognizable form or feature of the earth’s surface having a characteristic shape and produced by natural causes. Significance.—Geographic order suggests natural relationships. Running water, with weathering and gravitation, commonly sculptures landforms within a landscape. Over the ages, earthy material is removed from some landforms and deposited on others. Landforms are interrelated. An entire area has unity through the interrelationships of its landform. Each landform may have one kind of soil present or several. Climate, vegetation, and time of exposure to weathering of the parent materials are commonly about the same throughout the extent of the landform, depending on the relief of the area. Position on the landform may have influenced the soil-water relationships, microclimate, and vegetation. The proper identification of the landform is an important part of understanding the formative history of the soil and the materials from which the soil formed. This aids in the development of the soil mapping model and in the transfer of information between areas. Landform terms are also used as local phase criteria for separating components or uniquely naming soil map units. See Part 627 of this handbook for more information on naming physiographic phases. Classes.—The allowable list of landform terms is included in the NASIS data dictionary. Definitions of the terms are included in Part 629 of this handbook. Entries.—Enter the appropriate class name for the landforms on which each map unit component occurs. A representative value (term) may be indicated. It is possible to indicate the presence of one landform occurring on another landform (e.g., a dune on a flood plain). 618.41 Linear Extensibility Percent Definition.—“Linear extensibility percent” is the linear expression of the volume difference of natural soil fabric at 1/3-bar or 1/10-bar water content and oven dryness. The volume change is reported as percent change for the whole soil. Classes.—Shrink-swell classes are based on the change in length of an unconfined clod as moisture content is decreased from a moist to a dry state. If this change is expressed as a percent, the value used is LEP, linear extensibility percent. If it is expressed as a fraction, the value used is COLE, coefficient of linear extensibility. The shrink-swell classes are defined as follows: Shrink-Swell Class LEP COLE Low <3.0 <0.03 Moderate 3.0-5.9 0.03-0.06 High 6.0-8.9 0.06-0.09 Very High ≥9.0 ≥0.09 Significance.—If the shrink-swell classes is rated moderate to very high, shrinking and swelling can damage buildings, roads, and other structures. The high degree of shrinkage associated with high and very high shrink-swell classes can damage plant roots. Linear extensibility (expressed as cm of extension per meter of soil) is used for soil classification in the required characteristics for Vertic subgroups. Such soils will typically have LEP values of 6 or more. Measurement.— Coefficient of linear extensibility is measured directly as the change in clod dimension from moist to dry conditions and is expressed as a percentage of the volume change to the dry length: COLE = (moist length - dry length)/dry length When expressed as LEP (linear extensibility percent): LEP = COLE x 100 Linear extensibility may be determined by any of the following methods: For the core method of measurement, select a sample core from a wet or moist soil. Carefully measure the wet length of the cores and set the core upright in a dry place. If the core shrinks in a symmetrical shape without excessive cracking or crumbling, its length can be measured and linear extensibility percent calculated. If the core crumbles or cracks, measurements cannot be accurately determined by this method. In the coated clod method of measurement, shrink-swell potential can be estimated from the bulk density of soil measured when moist and when dry. The coated clod method is widely used and is the most versatile procedure for determining the bulk density of coherent soils. Procedures and calculations are given in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Linear extensibility percent can be calculated from bulk density moist (Dbm) and bulk density dry (Dbd) using the following formula: LEP = 100 [(Dbd/Dbm)1/3 -1] [1-(Volume % > 2 mm/100)] This equation is used to simplify the determination of shrink-swell potential classes. The classes are as follows: Dbd/Dbm Shrink-Swell Potential < 1.10 Low 1.10 - 1.20 Moderate 1.20 - 1.30 High ≥1.30 Very High Estimates.— Field estimates of shrink-swell potential can be made by observing desiccation cracks, slickensides, gilgai, soil creep, and leaning utility poles. Shrink-swell potential correlates closely with the kind and amount of clay. The greatest shrink-swell potential occurs in soils that have high amounts of 2:1 lattice clays, such as clay minerals in the smectite group. Illitic clays are intermediate, and kaolinitic clays are least affected by volume change as the content in moisture changes. Entries.—Enter the low, high, and representative linear extensibility percent values. If laboratory measurements or accurate field estimates are available, the high and low values do not need to correspond with shrink-swell class limits. However, if data is limited, the high and low values may correspond to the high and low limits of the appropriate shrink-swell class. The range of valid entries is from 0 to 30 percent, and tenths (one decimal place) are allowed. 618.42 Liquid Limit Definition.—“Liquid limit” is the water content of the soil material, which passes a no. 40 sieve, at the change between the liquid and plastic states. Significance.—The plasticity chart, given in ASTM Designation D 2487, is a plot of liquid limit (LL) versus plasticity index (PI) and is used in classifying soil in the Unified soil classification system. The liquid limit is also a criterion for classifying soil in the AASHTO classification system, as shown in Part 618, Subpart B, Exhibits, Section 618.83. Generally, the amount of clay- and silt-size particles, the organic matter content, and the type of minerals determine the liquid limit. Soils that have a high liquid limit have the capacity to hold a lot of water while maintaining a plastic or semisolid state. Measurement.— Tests are made on thoroughly puddled soil material that has passed a no. 40 (.425 mm) sieve and is expressed on a dry weight basis, according to ASTM Designation D 4318. This procedure requires the use of a liquid limit device, a special tool designed to standardize the arbitrary boundary between a liquid and plastic state of a soil. Estimates of liquid limit are made on soils during soil survey investigations and mapping activities. The liquid limit is usually inferred from clay mineralogy and clay content. If soils are tested later in the laboratory, adjustments are made to the field estimates as needed. Generally, experienced personnel can estimate these values with a reasonable degree of accuracy. Entries.—Enter the high, low, and representative values as a range for each horizon. The range of valid entries is from 0 to 400 percent, and tenths (one decimal place) are allowed. However, entries should be rounded to the nearest 10 percent unless they represent measured values. Enter “0” for nonplastic soils. The liquid limit for organic soil material is not defined and is assigned “null.” A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.102. 618.43 Organic Matter Definition.—“Organic matter percent” is the weight of decomposed plant and animal residue and expressed as a weight percentage of the soil material less than 2 mm in diameter. Significance Organic matter influences the physical and chemical properties of soils far more than the proportion to the small quantities present would suggest. The organic fraction influences plant growth through its influence on soil properties. It encourages granulation and good tilth, increases porosity and lowers bulk density, promotes water infiltration, reduces plasticity and cohesion, and increases the available water capacity. It has a high capacity to adsorb and exchange cations and is important to pesticide binding. It furnishes energy to micro-organisms in the soil. As it decomposes, it releases nitrogen, phosphorous, and sulfur. The distribution of organic carbon according to depth indicates different episodes of soil deposition or soil formation. Soils that are very high in organic matter have poor engineering properties and subside upon drying. Measurement.— Laboratory measurements are made using a dry combustion method to determine percent total carbon. For an estimate of organic carbon in calcareous soils, the carbon present in carbonate compounds, such as CaCO3, must be calculated and then subtracted from the total carbon. This is done using the equation: percent organic carbon = percent total carbon - [% <2 mm CaCO3 x 0.12]. The results are given as the percent of organic carbon in dry soil. To convert the figures for organic carbon to those for organic matter, multiply the organic carbon percentage by 1.724. The detailed procedures are outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Estimates.— Color and “feel” are the major properties used to estimate the amount of organic matter. Color comparisons in areas of similar materials can be made against laboratory data so that a soil scientist can make estimates. In general, black or dark colors indicate high amounts of organic matter. The contrast of color between the A horizon and subsurface horizons is also a good indicator. Entries.—Enter the high, low, and representative values for the range in organic matter in each horizon. The range of valid entries is from 0 to 100 percent and hundredths (two decimal places) are allowed. 618.44 Parent Material, Kind, Modifier, and Origin Parent material is the unconsolidated material, mineral or organic, from which the soil develops. The soil surveyor considers parent material in developing a model to be used for soil mapping. Soil scientists and specialists in other disciplines use parent material data to help interpret soil boundaries and project performance of the material below the soil. Many soil properties relate to parent material. Among these properties are proportions of sand, silt, and clay; chemical content; bulk density; structure; and the kinds and amounts of fragments. These properties affect interpretations and may be criteria used to separate soil series. Soil properties and landscape information infer parent material. Three data elements – parent material kind, parent material modifier, and parent material origin – describe parent material. Parent Material Kind Definition.—“Parent material kind” is a term describing the general physical, chemical, and mineralogical composition of the material, mineral or organic, from which the soil develops. Mode of deposition, weathering, or both may be implied or implicit. Classes.—The list of allowable entries is included in the NASIS data dictionary. Definitions of many of these terms are included in Part 629 of this handbook. Entries.—Enter the applicable class names for each map unit component. Multiple entries are permissible. Multiple rows of parent materials may also be indicated for a single component, such as loess over till over residuum. Parent Material Modifier Definition.—“Parent material modifier” is the general description of the texture of the parent material. Class limits correspond to those of the general texture defined in the Soil Survey Manual and the family category particle-size classes defined in Soil Taxonomy. Classes.—The classes of parent material modifiers are as follows: Clayey Coarse-loamy Coarse-silty Fine-loamy Fine-silty Gravelly Loamy Sandy Sandy and gravelly Sandy and silty Silty Silty and clayey Entries.—Enter the appropriate class name to modify the corresponding row of parent material kind. Parent Material Origin Definition.—“Parent material origin” is the type of bedrock from which the parent material is derived. Classes.—The allowable class names are included in the NASIS data dictionary and are the same as for the “bedrock kind” data element. Entries.—Enter the appropriate parent material origin class names that correspond to each parent material kind. Although this data element is intended to be used when “residuum” is the chosen parent material kind, it may also be used with other kinds of parent material. 618.45 Particle Density Definition.—“Particle density” is the mass per unit of volume of the solid soil particle, either mineral or organic. Also known as specific gravity. Significance.—Particle density is used in the calculation of weight and volume for soil (porosity). The relationship between bulk density, percent pore space, and the rate of sedimentation of solid particles in a liquid depends on particle density. The term “particle density” indicates wet particle density or specific gravity. Measurement.—The standard methods of measurement for particle density are: the ASTM standard test method for specific gravity of soils, ASTM Designation D 854-92, which uses soil materials passing a no. 4 sieve; the method described by Blake and Hartge in Methods of Soil Analysis, Part 1, Agronomy 9; or the method for volcanic soils described by Bielders and others in Soil Science Society of America Journal 54, pages 822-826. Estimates Particle density is often assumed to be 2.65 g cm-3; however, many minerals and material of various origins exhibit particle densities less than or greater than this standard. Particle density (Dp) may be calculated using the extractable iron and the organic carbon percentages in the following formula: OC is the organic carbon percentage and Fe is the percent extractable iron determined by dithionite-citrate extraction, or by an equivalent method. The particle density of the organic matter (Dp1) is assumed to be 1.4 g cm-3, that of the minerals from which the extractable iron originates (Dp2) is assumed to be 4.2 g cm-3, and that of the material exclusive of the organic matter and the minerals contributing to the extractable Fe (Dp3) is assumed to be 2.65 g cm-3. Entries.—Enter the representative value for the horizon. The range of valid entries is from 0.01 to 5 g cm-3, and hundredths (two decimal places) are allowed. 618.46 Particle Size Definition.—“Particle size” is the effective diameter of a particle as measured by sedimentation, sieving, or micrometric methods. Particle sizes are expressed as classes with specific effective diameter class limits. The broad classes are clay, silt, and sand, ranging from the smaller to the larger of the less than 2 mm mineral soil fraction. It includes fragments of weathered or poorly consolidated fragments that disperse to particles less than 2 mm. Significance.—The physical behavior of a soil is influenced by the size and percentage composition of the size classes. Particle size is important for most soil interpretations, for determination of soil hydrologic qualities, and for soil classification. Measurement.— Particle size is measured by sieving and sedimentation. The method used is method 3A1, which is outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.103. Classes.—The USDA uses the following size separates for the <2 mm mineral material: Definition.—“Total clay percentage” is the weight percentage of the mineral particles less than 0.002 mm in equivalent diameter in the less-than-2-mm soil fraction. Most of the material is in one of three groups of clay minerals or in a mixture of these clay minerals. The groups are kaolinite, smectite, and hydrous mica, the best known member of which is illite. Significance.—Physical and chemical activities of a soil are related to the kind and amount of clay minerals. Clay particles may have thousands of times more surface area per gram than silt particles and nearly a million times more surface area than very coarse sand particles. Thus, clay particles are the most chemically and physically active part of mineral soil. Clay mineralogy and clay percentage have a strong influence on engineering properties and the behavior of soil material when it is used as construction or foundation material. They influence linear extensibility, compressibility, bearing strength, and saturated hydraulic conductivity. The kind and amount of clay influence plant growth indirectly by affecting available water capacity, water intake rate, aeration, cation-exchange capacity, saturated hydraulic conductivity, erodibility, and workability. Up to a certain point, an increase in the amount of clay in the subsoil is desirable. Clay can increase the amount of water and nutrients stored in that zone. By slightly slowing the rate of water movement, it can reduce the rate of nutrient loss through leaching. If the amount of clay is great, it can impede water and air movement, restrict root penetration, increase runoff and, on sloping land, result in increased erosion. Clay particles are removed by percolating water from surface and subsurface horizons and deposited in the subsoil horizons. The amount of clay accumulation and its location in the profile provide clues for the soil scientist about soil genesis. Irregular clay distribution as related to depth may indicate lithologic discontinuities, especially if accompanied by irregular sand distribution. Measurement.— Clay content is measured in the laboratory by the pipette or hydrometer methods after the air-dry soil is pretreated to remove organic matter and soluble salts. Field estimates of clay content are made by manual methods. The way a wet soil ribbons (develops a long continuous ribbon) when pressed between the thumb and fingers gives a good idea of the amount of clay present. Excessive amounts of sodium can toughen the soil, making the soil feel more clayey. Care should be taken not to overestimate the amount of clay in sodic soils. Accuracy depends largely on frequent and attentive observation. Texture reference samples determined in the laboratory are used by soil scientists to calibrate the feel of soils with various percentages of clay. Entries.—Enter the high, low, and representative values of the clay total separate as a percent of the material less than 2 mm in size for each horizon. Enter a “0” if the amount is not significant, as in organic layers or in some andic soil materials. The range of valid entries is from 0 to 100 percent and tenths (one decimal place) are allowed. The representative value chosen should equate to a valid clay total content for the representative texture class posted for each horizon. Sand Percentage Definition.—“Sand percentage” is the weight percentage of the mineral particles less than 2 mm and greater than or equal to 0.05 mm in equivalent diameter in the less than 2 mm soil fraction. The sand separates recognized are very coarse, coarse, medium, fine, very fine, and total. Respective size limits are shown in Section 618.42(D) above. Much of the sand fraction is composed of fragments of rocks and primary minerals, especially quartz. Therefore, the sand fraction is quite chemically inactive. Significance.—Physical properties of the soil are influenced by the amounts of total sand and of the various sand fractions present in the soil. Sand particles, because of their size, have a direct impact on the porosity of the soil. This impact influences other properties, such as saturated hydraulic conductivity, available water capacity, water intake rates, aeration, and compressibility related to plant growth and engineering uses. Measurement.— Sand content is measured in the laboratory by the wet sieving method and then fractionated by dry sieving. Field estimates are made by manual methods. The degree of grittiness in a wet soil sample, when worked between the thumb and forefinger, gives an estimate of the sand content. The size of sand grains may be observed with the naked eye or with the aid of a hand lens. Entries.—Enter the high, low, and representative value of the sand total separate and each sand size separate (sand very coarse separate, sand coarse separate, sand medium separate, sand fine separate, and sand very fine separate) as a percent of the material less than 2 mm in size for each horizon. The sum of the representative values for the five sand size fractions must equal the representative value for the sand total separate. The range of valid entries is from 0 to 100 percent and tenths (one decimal place) are allowed. Enter a “0” if the amount is not significant, as in organic layers or in some andic soil materials. The representative values chosen should equate to a valid sand total content and sand size fraction content for the representative texture class posted for each horizon. Silt Percentage Definition.—“Silt percentage” is the weight percentage of the mineral particles greater than or equal to 0.002 mm but less than 0.05 mm in the less than 2 mm soil fraction. The silt separates recognized are fine, coarse, and total. The respective size limits are listed in paragraph 618.42(D) above. The silt separate is dominated by primary minerals, especially quartz, and therefore has a low chemical activity. Significance.—The silt separate possesses some plasticity, cohesiveness, and absorption, but to a much lesser degree than the clay separate. Silt particles act to slow water and air movement through the soil by filling voids between sand grains. A very high content of silt in a soil may be physically undesirable for some uses unless supplemented by adequate amounts of sand, clay, and organic matter. Measurement The silt content is measured in the laboratory in two phases. The fine silt is measured using the pipette method on the suspension remaining from the wet sieving process. Aliquots of the diluted suspension are removed at predetermined intervals based on Stokes Law. The aliquots are then dried and weighed. The coarse silt fraction is the difference between 100 percent and the sum of the sand, clay, and fine silt percentages. The silt content may be estimated in the field using the ribbon test as described for clay. The content of silt is usually estimated by first estimating the clay and sand portions and then subtracting that number from 100 percent. Silt tends to give the soil a smooth feel. Entries.—Enter the high, low, and representative value of the silt total separate and each silt size separate (silt coarse separate and silt fine separate) as a percent of the material less than 2 mm in size for each horizon. The sum of the representative values for the two silt size fractions must equal the representative value for the silt total separate. The range of valid entries is from 0 to 100 percent, and tenths (one decimal place) are allowed. Enter a “0” if the amount is not significant, as in organic layers or in some andic soil materials. The representative value chosen should equate to a valid silt total content for the representative texture class posted for each horizon. 618.47 Percent Passing Sieves Definition.—The percent passing sieve numbers 4, 10, 40, and 200 is the weight of material that passes through these sieves, based on the material less than 3 inches (75 mm) in size and expressed as a percentage. Significance.—Data for the percent passing sieves are used to classify the soil in the engineering classifications and to make judgments on soil properties and performance. Many soil characteristics are influenced by the depth distribution of grain sizes for the soil as well as the soil’s mode of deposition, stress history, density, and other features. Measurement.— Measurements involve sieve analysis for the determination of grain size distribution of that portion of the soil having particle diameters between 3 inches and 0.074 mm (no. 200 sieve). ASTM Designations D 422, C 136, and C 117 are applicable procedures. Measurements are made on most benchmark soils and other representative soils in survey areas. Estimates Estimates of the content of sand, silt, clay, and rock fragments that are made for soils during soil survey investigations and mapping activities are used to estimate percent passing sieves. If samples are tested later in a laboratory, adjustments are made to the field estimates as needed. Generally, experienced personnel can estimate these values with a high degree of accuracy. Estimates for percent passing sieves can be made from soil texture using the following general guidance: The percent passing #10 equals the less-than-2-mm fraction, and soil texture is based on the less than 2 mm fraction. Since sieves represent the less-than-3-inch fraction, the #40 and #200 sieve estimates must be adjusted when the percent passing #10 is less than 100 percent. The percent passing #40 and #200 that is determined by texture must be adjusted by multiplying the percent passing #40 and percent passing #200 by the percent passing #10. Pararock fragments are not cemented strongly enough to be retained on sieves. They are crushed and estimated into percent passing sieves. ASTM procedures use a roller crusher as a pretreatment of the soil material prior to sieving. Field estimates should try to replicate this procedure. Discrete artifacts which are either noncohesive or nonpersistent (e.g., paper) are not considered in estimating sieve values. Entries.—Enter the high, low, and representative values to represent the range of percent passing each sieve size for each horizon. The range includes the estimating accuracy as well as the range of values for a soil. The range of valid entries is from 0 to 100 percent, and tenths (one decimal place) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.104. 618.48 Plasticity Index Definition.—“Plasticity index” is the numerical difference between the liquid limit and the plastic limit. It is the range of water content in which a soil exhibits the characteristics of a plastic solid. The plastic limit is the water content that corresponds to an arbitrary limit between the plastic and semisolid states of a soil. Significance.—The plasticity index, when used in connection with the liquid limit, serves as a measure of the plasticity characteristics of a soil. The plasticity chart, given in ASTM Designation D 2487, is a plot of the liquid limit (LL) versus the plasticity index (PI) and is used in classifying soil in the Unified soil classification system. The plasticity index is also a criterion for classifying soil in the AASHTO classification system, as shown in Part 618, Subpart B, Exhibits, Section 618.83. Soils that have a high plasticity index have a wide range of moisture content in which the soil performs as a plastic material. Highly and moderately plastic clays have large PI values. Measurements.—Tests are made on that portion of the soil having particles passing the no. 40, (425 micrometer) sieve, according to ASTM Designation D 423. Measurements are made on most benchmark soils and on other representative soils in survey areas. Estimates of plasticity index are made on all soils during soil survey investigations and mapping activities. The plasticity index is usually not estimated directly: a position on the plasticity chart in ASTM Designation D 2487 is estimated and the plasticity index is determined from the chart. If soils are later tested in the laboratory, adjustments are made to the field procedures as needed. Generally, experienced personnel can estimate these values with a reasonable degree of accuracy. Estimates are expressed in ranges that include the estimating accuracy as well as the range of values from one location to another within the map unit. Entries.—Enter the high, low, and representative values to represent the range of plasticity index for each horizon. The range of valid entries is from 0 to 130 percent, and tenths (one decimal place) are allowed. However, entries should be rounded to the nearest 5 percent unless they represent measured values. Enter “0” for nonplastic soils. The plasticity index for organic soil material is not defined and is assigned “null.” A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.102. 618.49 Ponding Depth, Duration Class, Frequency Class, and Month Ponding is standing water in a closed depression. The water is removed only by deep percolation, transpiration, evaporation, or by a combination of these processes. Ponding of soils is classified according to depth, frequency, duration, and the beginning and ending months in which standing water is observed. Ponding Depth Definition.—“Ponding depth” is the depth of the surface water that is ponding on the soil. Entries.—Enter the high, low, and representative values for the ponding depth, in centimeters, for the map unit component. The range of valid entries is from 0 to 185 cm, and only whole numbers (integers) are allowed. Ponding Duration Class Definition.—“Ponding duration class” is the average duration, or length of time, of the ponding occurrence. Classes.—The ponding duration classes are: Ponding Duration Class Duration of the Ponding Occurrence Very brief 4 to 48 hours Brief 2 to 7 days Long 7 to 30 days Very long > 30 days Entries.—Enter very brief, brief, long, or very long for the map unit component. Only use entries if ponding frequency (defined below) occurs more often than rare. Ponding Frequency Class Definition.—“Ponding frequency class” is the number of times ponding occurs over a period of time. Classes.—The ponding frequency classes are: Ponding Frequency Class Definition None No reasonable possibility of ponding; near 0 percent chance of ponding in any year. Rare Ponding is unlikely but possible under unusual weather conditions; nearly 0 to 5 percent chance of ponding in any year or nearly 0 to 5 times in 100 years. Occasional Ponding is expected infrequently under usual weather conditions; 5 to 50 percent chance of ponding in any year or nearly 5 to 50 times in 100 years. Frequent Ponding is likely to occur under usual weather conditions; more than a 50 percent chance in any year (i.e., 50 times in 100 years). Entries.—Enter none, rare, occasional, or frequent as appropriate for the map unit component. Ponding Month Definition.—“Ponding month” is the calendar months in which ponding is expected. Classes.—The time of year when ponding is likely to occur is expressed in months for the expected beginning to expected end of the ponding period. The time period expressed includes two-thirds to three-fourths of the occurrences. Entries.—Yearly ponding frequency classes are assigned to months and indicate the months of occurrence and not the frequency of the ponding during the month. Enter the name of each month of the year in which ponding is expected. Significance.—The susceptibility of soils to ponding is important for homes, building sites, and sanitary facilities. Time and duration of the ponding are critical factors in determining plant species. Ponding during the dormant season has few if any harmful effects on plant growth or mortality and may even improve growth. Estimates.— Generally, estimates of ponding frequency and duration can be made for each soil. Where the natural infiltration, saturated hydraulic conductivity, and surface and subsurface drainage of soils is altered, ponding studies are needed to reflect present ponding characteristics. Evidence of ponding events should be gathered during soil survey fieldwork. High water lines and other signs of maximum water height are recorded. Other records may also exist. Certain landform features are subject to ponding. These features are characteristics of closed drainage systems and include potholes, playas, sloughs, and backswamps. Most of these features are recognizable when correlating features on aerial photographs with ground observations. Different kinds of vegetation and soils are normally associated with these geomorphic features. The vegetation that grows in ponded areas may furnish clues to past ponding and indicate the potential for ponding in the future. Generally, native vegetation in ponded areas consists of obligate and facultative wet hydrophytes. Some plant species are intolerant of ponding and do not grow in areas that are ponded. The soil provides clues to past ponding, but characteristics vary according to climate and soil conditions. Some of the clues (alone or in any combination) are— A dark surface horizon or layer overlying a gleyed subsoil. Many prominent redoximorphic features that have low value and chroma. Capillary transport and concentrations of carbonates or sulfates, or both, in the upper soil horizons. Dark colors and high levels of organic matter throughout the profile. 618.50 Pores “Pore space” is a general term for voids in the soil material. The term includes matrix, nonmatrix, and interstructural pore space. For water movement at low suction and conditions of satiation, the nonmatrix and interstructural porosity have particular importance. Matrix Pores.—Matrix pores are formed by the agents that control the packing of the primary soil particles (i.e., primary packing voids). These pores are usually smaller than nonmatrix pores. Additionally, their aggregate volume and size can change markedly according to water state for soil horizons or layers with high extensibility. Nonmatrix Pores.—Nonmatrix pores are relatively large voids that are expected to be present both when the soil is moderately moist or wetter as well as in drier states. The voids are not bounded by the planes that delimit structural units. Nonmatrix pores may be formed by roots, animals, the action of compressed air, and other agents. The size of the distribution of nonmatrix pores usually bears no relationship to the particle-size distribution and the related matrix pore-size distribution. Interstructural Pores.—Interstructural pores are delimited by structural units. Inferences as to the interstructural porosity may be obtained from the structure description. Commonly, interstructural pores are at least crudely planar. Nonmatrix pores are described by quantity, size, shape, and vertical continuity (generally in that order). Pore Quantity Definition.—“Pore quantity” is defined by classes that pertain to the number of a selected size of pores per unit area of undisturbed soils. The unit area that is evaluated varies according to the size class of the pores: 1 cm2 for very fine and fine pores, 1 dm2 for medium and coarse pores, and 1 m2 for very coarse pores. Classes.—The pore quantity classes are: Pore Quantity Class Number of Pores per Unit Area Few < 1 Common ≥ 1-5 Many ≥ 5 Entries.—Enter pore quantity as pores/area. Enter the high, low, and representative values as whole numbers between 0 and 99 for the horizon. Pore Size Definition.—“Pore size” is the average diameter of the pore. Classes.—The pore size classes are: Pore Size Class Pore Size (mm) Very fine < 1 Fine 1 - <2 Medium 2 - <5 Coarse 5 - <10 Very coarse ≥ 10 Entries.—Enter a single class or classes for the horizon. Pore Shape Definition.—“Pore shape” is a description of the multiarial shape of the pore. The shapes of nonmatrix pores are dendritic tubular (approximately cylindrical, elongated, and branching), irregular (nonconnected cavities or chambers), tubular (approximately cylindrical and elongated), or vesicular (approximately spherical or elliptical). The primary packing voids between soil particles or rock fragments are referred to as interstitial pores. Classes.—The pore shape classes are: Dendritic tubular Interstitial Irregular Tubular Vesicular Entries.—Enter one of the classes from the pore shape list for the horizon. Vertical Continuity Definition.—“Vertical continuity” is the average vertical distance through which the minimum pore diameter exceeds 0.5 mm when the soil layer is moist or wetter. Classes.—The vertical continuity classes are: Vertical Continuity Class Vertical Distance (cm) Low < 1 Moderate 1 - <10 High ≥ 10 Entries.—Enter one of the vertical continuity classes. 618.51 Reaction, Soil (pH) Definition.—“Soil reaction” is a numerical expression of the relative acidity or alkalinity of a soil. Classes.—The descriptive terms for reaction and their respective ranges in pH are: Reaction Class Range in pH Ultra acid 1.8-3.4 Extremely acid 3.5-4.4 Very strongly acid 4.5-5.0 Strongly acid 5.1-5.5 Moderately acid 5.6-6.0 Slightly acid 6.1-6.5 Neutral 6.6-7.3 Slightly alkaline 7.4-7.8 Moderately alkaline 7.9-8.4 Strongly alkaline 8.5-9.0 Very strongly alkaline 9.1-11.0 Significance A principal value of soil pH is the information it provides about associated soil characteristics. Two examples are phosphorus availability and base saturation. Soils that have a pH of approximately 6 or 7 generally have the most ready availability of plant nutrients. Strongly acid or more acid soils have low extractable calcium and magnesium; a high solubility of aluminum, iron, and boron, and a low solubility of molybdenum. In addition, these soils may possibly have organic toxins and generally have a low availability of nitrogen and phosphorus. At the other extreme are alkaline soils. Calcium, magnesium, and molybdenum are abundant where there is little or no toxic aluminum and nitrogen is readily available. If pH is above 7.9, the soils may have an inadequate availability of iron, manganese, copper, zinc, and especially phosphorus and boron. Soil reaction is one of several properties used as a general indicator of soil corrosivity or the soil’s susceptibility to dispersion. In general, soils that are either highly alkaline or highly acid are likely to be corrosive to steel. Soils that have pH <5.5 are likely to be corrosive to concrete. Soils that have pH >8.5 are likely to be highly dispersible and may have a piping problem. Soil reaction is used for soil classification in the required characteristics for sulfidic materials, in the key to calcareous and reaction classes for mineral soils, in the key to reaction classes for Histosols and Histels, and in criteria for certain taxa such as Sulfic subgroups. Measurement.—The most common soil laboratory measurement of pH is the 1:1 water method. In this method, a crushed and sieved soil sample is mixed with an equal amount of water and a measurement is made of the suspension using a pH meter. Another common method, used for mineral and organic soils, is the 0.01M calcium chloride method. A new method to indicate the possible presence of sulfidic materials is the hydrogen peroxide test, delta pH for acid sulfate soils. This method uses hydrogen peroxide to rapidly oxidize sulphur compounds which releases elemental sulphur and quickly decreases the pH. In NASIS, the pH values derived from these three methods are populated in separate data elements. The pH values derived from water suspension are affected by field applications of fertilizer or other salts in the soil, the content of carbon dioxide in the soil, and the moisture content at the time of sampling. The 0.01M calcium chloride method reduces these influences. The laboratory procedure for measuring pH by the 1:1 water and 0.01 M calcium chloride methods are described in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Estimates.— A variety of field test kits are available for determination of pH in the field. The methods include a water-soluble dye, which is mixed with soil and thus produces a color that is compared with a chart; a dye-impregnated paper, which changes color according to differences in pH; and portable glass electrodes. Each State office can recommend a suitable pH method for the soils in the State. If requested, the NSSC Kellogg Soil Survey Laboratory makes suggestions for suitable methods for field measurements and furnishes NRCS soil scientists with the proper chemicals. Entries.—Soil reaction (pH) is time and moisture dependent, and water pH can vary up to a whole unit during the growing season. The range of pH should reflect the variations. The 1:1 water method generally is used with mineral soils. Mineral and organic soils are measured in a 1:2 0.01M calcium chloride solution and with the hydrogen peroxide test, delta pH method. Separate entries are made by horizon for pH 1:1 water, pH 1:2 0.01M calcium chloride, and final pH oxidized, as needed. Enter the high, low, and representative values of the appropriate estimated pH range for each horizon. If laboratory measurements or accurate field estimates are available, the high and low values do not need to correspond with reaction class limits. However, if data is limited, then pH values may reflect reaction class limits such as 1.8-3.4, 3.5-4.4, etc., or a combination of reaction classes, such as 4.5-5.5, can be entered. 618.52 Restriction Kind, Depth, Thickness, and Hardness Restriction Kind Definition.—“Restriction kind” is the type of nearly continuous layer that has one or more physical, chemical, or thermal properties that significantly reduce the movement of water and air through the soil or that otherwise provide an unfavorable root environment. Bedrock (e.g., limestone), cemented horizons (e.g., duripan), densic material (e.g., dense till), frozen horizons or layers (e.g., permanent ground ice), and horizontally oriented, human-manufactured materials (e.g., concrete) are examples of subsurface layers that are kinds of restrictions. Significance.—Restrictive layers limit plant growth by restricting the limits of the rooting zone. They also impede or restrict the movement of soil water vertically through the soil profile and have a direct impact on the quality and quantity of ground water and surface water. Restrictions are important for both soil interpretations and soil classification. Measurement.—Identify and describe restrictive soil layers in the field. Observe, measure, and record the restriction kind along with their depth, thickness, and hardness (defined below). When describing pedons, identify types or kinds of restrictions by suffix symbols, such as "d," "f," "m," "r," "v," or "x," or by the master layers “M” or "R." Use measurements or observations made throughout the extent of occurrence of a soil as a basis for estimates of restriction kind. Entries.—Enter the appropriate choice for the kind of restrictive horizon or layer from the following list— Abrupt textural change Bedrock, densic Bedrock, lithic Bedrock, paralithic Cemented horizon Densic material Duripan Fragipan Human-manufactured materials Natric Ortstein Permafrost Petrocalcic Petroferric Petrogypsic Placic Plinthite Salic Strongly contrasting textural stratification Sulfuric Restriction Depth Definition.—“Restriction depth” is the vertical distance from the soil surface to the upper and lower boundaries of the restriction. Measurement.—Use measurements or observations made throughout the extent of occurrence of a soil as a basis for estimates of restriction depth. Entries.—Restriction depth values used to populate component data in NASIS are not specific to any one point. They are a reflection of commonly observed values based on field observations and are intended to model the component as it occurs throughout the map unit. Enter the high, low, and representative values for the top and bottom restriction depths in centimeters using whole numbers (integers). Restriction Thickness Definition.—“Restriction thickness” is the distance from the top to the bottom of a restrictive layer. Significance.—Restriction thickness has a significant impact on the ease of mechanical excavation. Measurement.—Use observations made throughout the extent of occurrence of a soil as a basis for estimates of restriction thickness. Entries.—Restriction thickness values used to populate component data in NASIS are not specific to any one point. They are a reflection of commonly observed values based on field observations and are intended to model the component as it occurs throughout the map unit. Enter the high, low, and representative values for the thickness in centimeters. The range of valid entries is from 1 to 999, and only whole numbers (integers) are allowed. Restriction Hardness Definition.—“Restriction hardness” is the rupture resistance cemented of an air-dried, then submerged block-like specimen of mineral material. Ice is not applicable. Significance.—Restriction hardness has a significant impact on the ease of mechanical excavation. Use excavation difficulty classes (defined above) to evaluate the relationships of restriction layers to excavations. Measurement.—Use observations made throughout the extent of occurrence of a soil as a basis for estimates of restriction hardness. For measurements of the restriction hardness, use the procedures and classes of cementation listed with the rupture resistance classes. Classes are described for like specimens about 25-30 mm on edge that are air-dried and then submerged in water for at least 1 hour. Compress the specimen between extended thumb and forefinger, between both hands, or between the foot and a nonresilient flat surface. If the specimen resists compression, drop a weight onto it from progressively greater heights until it ruptures. Failure is the point of the initial detection of deformation or rupture. Stress applied in the hand should be over a 1-second period. Learn the tactile sense of the class limits by applying force to top-loading scales and sensing the pressure through the tips of the fingers or through the ball of the foot. Use postal scales for the resistance range that is testable with the fingers. Use a bathroom scale for the higher rupture resistance range. Classes.—Restriction hardness is rated using the following classes and operation descriptions: Restriction Hardness (Rupture Resistance) Class Operation Description Noncemented Fails under very slight force applied slowly between thumb and forefinger (<8N). Extremely weakly cemented Fails under slight force applied slowly between thumb and forefinger (8 to 20N). Very weakly cemented Fails under moderate force applied slowly between thumb and forefinger (20 to 40N). Weakly cemented Fails under strong force applied slowly between thumb and forefinger (about 80N maximum force can be applied) (40 to 80N). Moderately cemented Cannot be failed between thumb and forefinger but can be failed between both hands or by placing specimen on a nonresilient surface and applying gentle force underfoot (80 to 160N). Strongly cemented Cannot be failed in hands but can be failed underfoot by full body weight (about 800N) applied slowly (160 to 800N). Very strongly cemented Cannot be failed underfoot by full body weight but can be failed by <3J blow (800N to 3J). Indurated Cannot be failed by blow of 3J (≥ 3J). Both force (Newtons, N) and energy (joules, J) are employed. The number of Newtons is 10 times the kilograms of force. One joule is the energy delivered by dropping a 1 kg weight a distance of 10 cm. 618.53 Saturated Hydraulic Conductivity Definition.—“Saturated hydraulic conductivity” is the ease with which pores of a saturated soil transmit water. Formally, it is the proportionality coefficient that expresses the relationship of the rate of water movement to hydraulic gradient in Darcy's Law (a law that describes the rate of water movement through porous media). It is expressed in micrometers per second. To convert micrometers per second to inches per hour, multiply micrometers per second by 0.1417. The historical definition of “saturated hydraulic conductivity” is the amount of water that would move vertically through a unit area of saturated soil in unit time under unit hydraulic gradient. Significance.—Saturated hydraulic conductivity is used in soil interpretations. It is also known as Ksat. Saturated hydraulic conductivity is used for soil classification in criteria for certain taxa such as the Albaqualfs and Albaquults great groups. Measurement.—Means of measurement, such as the Amoozemeter and double ring infiltrometers, provide some basis for estimation of saturated hydraulic conductivity. No method has been accepted as a standard. Since measurements are difficult to make and are only available for relatively few soils, estimates of saturated hydraulic conductivity are based on soil properties. Estimates.— The soil properties that affect saturated hydraulic conductivity are distribution, continuity, size, and shape of pores. Since the pore geometry of a soil is not readily observable or measurable, observable properties related to pore geometry are used to make estimates of saturated hydraulic conductivity. These properties are texture, structure, pore size, density, organic matter content, and mineralogy. Part 618, Subpart B, Exhibits, Section 618.88 provides a guide for estimating saturated hydraulic conductivity according to soil texture and bulk density or according to specified overriding conditions. In making estimates, the soil characteristic that exerts the greatest control for many soils is texture. The general relationships shown in Part 618, Subpart B, Exhibits, Section 618.88 are adjusted up or down depending on bulk density. Structure, pore size, organic matter content, clay mineralogy, and other features observed within the soil profile, such as consistency, dry layers in wet seasons, root mats or absence of roots, and evidence of perched water levels or standing water, are good field indicators for adjusting estimates. Measurement.— Slope aspect is measured clockwise from true north as an angle between 0 and 360 degrees. Tools such as geographic information systems (GIS) can be used to consistently predict and identify slope aspect. Entries.—For map unit components that are aspect dependent, enter the slope aspect counterclockwise, slope aspect clockwise, and slope aspect representative. The range of valid entries is from a minimum of 0 degrees to a maximum of 360 degrees. Record values to the nearest whole number (integer). The fields may be left NULL for those components that are not aspect dependent. “Slope aspect counterclockwise” is one end of the range in characteristics for the slope aspect of a component. This end of the range is expressed in degrees measured clockwise from true north, but in the direction counterclockwise from the representative slope aspect. “Slope aspect clockwise” is one end of the range in characteristics for the slope aspect of a component. This end of the range is expressed in degrees measure clockwise from true north, and in the direction clockwise from the representative slope aspect. “Slope aspect representative” is the common, typical, or expected direction toward which the surface of the soil faces, measured in degrees clockwise from true north. 618.55 Slope Gradient Definition.—“Slope gradient” is the difference in elevation between two points and is expressed as a percentage of the distance between those points. For example, a difference in elevation of 1 meter over a horizontal distance of 100 meters is a slope of 1 percent. Significance.—Slope gradient influences the retention and movement of water, the potential for soil slippage and accelerated erosion, the ease with which machinery can be used, soil-water states, and the engineering uses of the soil. Slope is used for soil classification in criteria for certain taxa such as the Fluvents suborder, Fluvaquents great group, Fluvaquentic and Fluventic subgroups, and several Cumulic subgroups. Measurement.— Slope gradient is usually measured with a hand level or clinometer. The range is determined by summarizing data from several sightings. Entries.—Enter the high, low, and representative values to represent the range of slope gradient as a percentage for the map unit component. The range of valid entries is from 0 to 999 percent, and tenths (one decimal place) are allowed but should only be used for values less than 1 percent. 618.56 Slope Length, USLE Definition.—“Slope length” is the horizontal distance from the origin of overland flow to the point where either the slope gradient decreases enough that deposition begins or runoff becomes concentrated in a defined channel. Refer to Agriculture Handbook 703. Significance.—Slope length has considerable control over runoff and potential accelerated water erosion. Slope length is combined with slope gradient in erosion prediction equations to account for the effect of topography on erosion. Measurement Slope length is measured from the point of origin of overland flow to the point where the slope gradient decreases enough that deposition begins or runoff becomes concentrated in a defined channel. In cropland, defined channels are usually ephemeral gullies or, in rare instances where they are near a field edge, are a classic gully or stream. Surface runoff will usually concentrate in less than 400 feet (120 meters), although longer slope lengths of up to 1000 feet are occasionally found. The maximum distance allowed in erosion equations is 1000 feet (305 meters). Conversion to the horizontal distance is made in the conversion process within the equation model. Assume no support practices. Ignore practices such as terraces or diversions. Slope length is best determined by pacing or measuring in the field. Do not use contour maps to estimate slope lengths unless contour intervals are 1 foot or less. Slope lengths estimated from contour maps are usually too long because most maps do not have the detail needed to indicate all ephemeral gullies and concentrated flow areas that end the slope lengths. Refer to figures 4-1 through 4-10 within Agriculture Handbook 703 for more landscape guidance. Entries.—Enter the high, low, and representative values for the range for each map unit component. Enter a whole number that represents the slope length in meters, from the point of origin of overland flow to the point of deposition or concentrated flow, of the slope on which the component lies. The slope length may be fully encompassed within one map unit or may cross several map units. The minimum value is 0, and the maximum value used in erosion equations is 305 meters. The NASIS database allows valid entries from 0 to 4000 meters. 618.57 Sodium Adsorption Ratio Definition.—“Sodium adsorption ratio” (SAR) is a measure of the amount of sodium (Na+) relative to calcium (Ca2+) and magnesium (Mg2+) in the water extracted from a saturated soil paste. It is the ratio of the Na concentration divided by the square root of one-half of the Ca + Mg concentration. SAR is calculated from the equation: SAR = Na+ / [(Ca2+ + Mg2+)/2]0.5 Significance.—Sodium adsorption ratio is used for soil classification in the required characteristics for the natric horizon, in the key to soil orders and key to suborders of Inceptisols and Mollisols, and in criteria for certain taxa such as Sodic subgroups. Soils that have values for sodium adsorption ratio of 13 or more may have an increased dispersion of organic matter and clay particles, reduced saturated hydraulic conductivity and aeration, and a general degradation of soil structure. Measurement.— The concentration of Na, Ca, and Mg ions is measured in a water extract from a saturated soil paste. The method is described in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. The sodium adsorption ratio is then calculated from the molar concentrations of the three cations using the equation shown above in Section 618.57 (A). Entries.—Enter the high, low, and representative values for the range of sodium adsorption ratio for each horizon. Enter “0” where SAR is negligible. The range of valid entries is from 0 to 9999, and tenths (one decimal place) are allowed. 618.58 Soil Erodibility Factors, USLE, RUSLE2 Definition.—Soil erodibility factors (Kw) and (Kf) quantify soil detachment by runoff and raindrop impact. These erodibility factors are indexes used to predict the long-term average soil loss from sheet and rill erosion under crop systems and conservation techniques. Factor Kw applies to the whole soil and factor Kf applies only to the fine-earth (less than 2.0 mm) fraction. The procedure for determining the Kf factor is outlined in Agriculture Handbook 703, Predicting Soil Erosion by Water: A Guide to Conservation Planning With the Revised Universal Soil Loss Equation (RUSLE), USDA, Agricultural Research Service, 1997. The K factors for soils in Hawaii and the Pacific Basin were extrapolated from local research. The nomograph, shown in Part 618, Subpart B, Exhibits, Section 618.91, was not used to determine K factors for soils in Hawaii. Significance.—Soil erodibility factors Kw or Kf are used in the erosion prediction equations USLE and RUSLE. Soil properties that influence rainfall erosion are those that affect— Infiltration rate, movement of water through the soil, and water storage capacity. Dispersion, detachability, abrasion, and mobility by rainfall and runoff. Some of the most important properties are texture, organic matter content, structure size class, and the saturated hydraulic conductivity of the subsoil. Estimates The Kw factor is measured by applying a series of simulated rainstorms on freshly tilled plots. Direct measurement of K factors is both costly and time consuming and is conducted only for a few selected soils. Reliable estimates of Kf factor are obtained from the soil erodibility nomograph, which is presented on page 11 of Agriculture Handbook 537 and reproduced in Part 618, Subpart B, Exhibits, Section 618.91, or by using the soil erodibility equation. The nomograph integrates the relationship between the Kf factor and the following five soil properties: Percent silt plus very fine sand Percent sand greater than 0.10 mm Organic matter content Soil structure Saturated hydraulic conductivity The soil erodibility equation which follows also provides an estimate of Kf. K factor = {2.1 x M1.14 X 10-4 x (12-a)+3.25 x (b-2)+2.5 x (c-3)}/100 where: M = (percent si + percent vfs) X (100 - percent clay). Example: For a soil with 29.0% silt, 12.3% very fine sand, and 36% clay M = (29.0+12.3) x (100-36) = 2,643.20. a = percent organic matter (0, 1, 2, 3, or 4). Use worse case organic matter content assuming long-term cultivation. b = soil structure code (1, = very fine granular, 2, = fine granular, 3, = med or coarse granular, or 4 = blocky, platy, or massive) c = profile saturated hydraulic conductivity code (1, 2, 3, 4, 5, or 6). Use the layer with the lowest Ksat rv (representative value) in the permeability control section. The permeability control section is the zone from the top of the mineral soil layer being evaluated to a depth of 50 cm below the top of that soil layer but should not exceed a profile depth of 200 cm. The permeability control section guarantees that a specific zone is only considered relative to the mineral soil layer being evaluated. Include the permeability of any bedrock or other non-soil layers in the permeability control section. Note that the codes were initially established using the 1951 Soil Survey Manual. The codes correspond to the following saturated hydraulic conductivity ranges: Profile permeability class code Permeability class of 1951 Saturated hydraulic conductivity range μm/sec Saturated hydraulic conductivity classes 1993 6 Very slow <0.30 very low or mod. low 5 Slow 0.30 to <1.20 mod. low 4 Slow or mod. 1.20 to <4.80 mod. high 3 Moderate 4.80 to <15.00 mod. high or high 2 Mod. or rapid 15.00 to <30.00 high 1 Rapid ≥30.00 high or very high The accuracy of the nomograph and equation has been demonstrated for a large number of soils in the United States. However, the nomograph and the equation may not be applicable to some soils having properties that are uniquely different from those used in developing the nomograph. For example, the nomograph does not accurately predict Kf factors for certain Oxisols in Puerto Rico or the Hawaiian Islands; some soils with andic properties, organic soil materials, or low activity clays; and some calcareous or micaceous soils. In these cases, Kf factors are estimated using the best information at hand and knowledge of the potential for rainfall erosion. See Agriculture Handbook 703 for more information. When using the nomograph and the equation, care should be taken to select an organic matter percentage that is most representative of the horizon being considered, assuming long-term cultivation. It is acceptable to use linear interpolations between plotted lines on the nomograph and values in hundredths (two decimal places) for organic matter content in the equation. For horizons that have organic matter content greater than 4 percent, use the 4 percent curve in the nomograph and exactly 4 percent in the equation. Rock or pararock fragments are not taken into account in the nomograph or the soil erodibility equation. If fragments are substantial, they have an armoring effect. Pararock fragments are assumed to break down with cultivation or other manipulation and so are not used in determining Kw factors. If a soil has mixtures of rock and pararock fragments, the Kw factor should reflect the degree of protection afforded only by the rock fragments. Guidelines for determining Kw factors are as follows: Then use the table in Part 618, Subpart B, Exhibits, Section 618.92 to convert the Kf value of the soil fraction less than 2 mm in diameter, which is derived from either the nomograph in Part 618, Subpart B, Exhibits, Section 618.91 or from the soil erodibility equation, to a Kw factor adjusted for the total volume of rock fragments. The Kw factor is adjusted only when the total content of rock fragment values in the layer, by volume, is equal to or greater than 15 percent. If total rock fragment content, by volume, is less than 15 percent, the Kw factor equals the Kf factor. In practice, the representative values (rv’s) for rock fragment volume, as populated in the NASIS Horizon Fragments Table, are summed for each size fraction to compute the total rock fragment content for the layer. If the soil on site contains more or less rock fragments than the mean of the range reported, adjustments can be made in Kf by using Part 618, Subpart B, Exhibits, Section 618.92. Select the estimates of total rock fragment volume percentages, and then use Part 618, Subpart B, Exhibits, Section 618.92. Enter Part 618, Subpart B, Exhibits, Section 618.92 in line with the rock fragment volume percentage and find, in the appropriate line, the nearest value to the Kf factor. Within that column, read the Kw factor on the line with the percentage of rock fragments of the soil for which you are making the estimate. Round the K factor displayed in the table to the closest acceptable K factor class entry, as shown below. This is the new Kw factor adjusted for rock fragments on site. The acceptable entries for Kw and Kf classes are 0.02, 0.05, 0.10, 0.15, 0.17, 0.20, 0.24, 0.28, 0.32, 0.37, 0.43, 0.49, 0.55, and 0.64. Use the comparison reports and calculation script in NASIS for help in populating Kf and Kw factors. Use the reports to print or export the currently stored values for Kf and Kw factor classes for each component in a selected set for comparison with the computed Kf and Kw factor classes using the calculation script formulas. The comparison reports give a preview of the results of the K factor calculations and should be used before the decision is made to run the calculation and save the new data. Soil horizons that do not have rock fragments are assigned equal Kw and Kf factors. In horizons where total rock fragments are 15 percent or more, by volume, the Kw factor is always less than the Kf factor. For example: Depth (in) USDA Texture Kw Kf 0-5 GR-L 0.20 0.32 0-5 L 0.32 0.32 0-5 GRV-L 0.10 0.32 0-46 CL 0.28 0.28 46-60 SL 0.20 0.20 Soils that have similar properties and erosivity should be grouped in similar K factor classes. 618.59 Soil Erodibility Factors for WEPP Soil erodibility factors for WEPP include interrill erodibility (Ki), rill erodibility (Kr), and critical hydraulic shear stress (Tc). These erodibility factors for the WEPP erosion model quantify the susceptibility of soil detachment by water. These erodibility factors predict the long-term average soil loss which results from sheet and rill erosion under various alternative combinations of crop systems and conservation techniques. The Ki, Kr, and Tc factors are used in a continuous simulation computer model which predicts soil loss and deposition on a hillslope. Reference NSERL Report No. 9, USDA, Agricultural Research Service, National Erosion Research Laboratory, August 1994, documentation version 94.7. This procedure does not include data for soils with highly weathered material (e.g., oxic horizons) and those with andic soil properties. These factors are quantitative and calculated using experimental equations. They are different than the soil erodibility factors used in USLE and RUSLE. Interrill erodibility (Ki) Definition.—“Interrill erodibility (Ki)” is the susceptibility of detachment and transport of soil particles by water. It is the susceptibility of the soil to movement to a rill carrying runoff. Significance.—Interrill erodibility (Ki) is a measure of sediment delivery rate to rills as a function of rainfall intensity. The Ki values for soil need to be adjusted if factors that influence the resistance of soil to detachment occur. These factors include live and dead root biomass, soil freezing and thawing, and mechanical and livestock compaction. Measurement.— Interrill erodibility (Ki) measurements are determined from rainfall simulation experiments. These experiments require the use of specialized equipment and specialized measurement techniques in a research setting. Calculations.—Use the following equations: For cropland soils with 30 percent or more sand: Ki = 2,728,000 + 192,100 x (% very fine sand) Where: Very fine sand must be less than or equal to 40 percent; if very fine sand is greater, use 40 percent. Definition.—“Rill erodibility (Kr)” is a measure of the susceptibility of a soil to detachment by flowing water. As rill erodibility (Kr) increases, rill erosion rates increase. Significance.—Rill erodibility (Kr) is often defined as the soil detachment per unit increase in shear stress of clear water flow. The rate of soil detachment in rills varies because of a number of factors, including soil disturbance by tillage, living root biomass, incorporated residue, fragments, soil consolidation, freezing and thawing, and wheel and livestock compaction. Measurement.— Rill erodibility (Kr) measurements are determined by rainfall simulation and flow simulation experiments. These experiments require the use of specialized equipment and specialized measurement techniques in a research setting. Calculations.—Use the following equations: For cropland soils with 30 percent or more sand: Kr = 0.00197 + 0.00030 x (% very fine sand) + 0.03863 x EXP(-1.84 x ORGMAT) Where: Organic matter (ORGMAT) is the organic matter in the surface soil (assuming that organic matter equals 1.724 times organic carbon content). Organic matter must exceed 0.35 percent; if less, use 0.35 percent. Very fine sand must be less than or equal to 40 percent; if greater, use 40 percent. Entries.—The computer generates the value by using the above formulas. Allowable Kr values range from 0.002 to 0.045 s/m. Critical shear stress (Tc) Definition.—“Critical shear stress (Tc)” is the hydraulic shear that must be exceeded before rill erosion can occur. Significance. Critical shear stress (Tc) is important in the rill detachment equation. It is the shear stress below which no soil detachment occurs. Critical shear stress (Tc) is the shear intercept on a plot of detachment by clear water versus shear stress in rills. Measurements.—Critical shear stress (Tc) is derived from a specialized research project. Calculations.—Use the following equations: For cropland soils with 30 percent or more sand: Tc = 2.67 + 0.065 x (% clay) - 0.058 x (% very fine sand) Where: Very fine sand must be less than or equal to 40 percent; if greater, use 40 percent. For cropland soils with less than 30 percent sand: Tc = 3.5 Entries.—No manual entry is needed. The value is computer generated using the above formulas. Allowable Tc values range from 1 to 6 N/m2. 618.60 Soil Moisture Status Definition.—“Soil moisture status” is the mean monthly soil water state at a specified depth. Classes.—The water state classes used in soil moisture status are dry, moist, and wet. These classes are defined as follows: Water State Class Definition Dry ≥15 bar suction Moist <15 bar to ≥0.0 bar (moist plus nonsatiated wet) Wet <0.0 bar; free water present (satiated wet) Significance.—Soil moisture status is a recording of the generalized water states for a soil component. Soil moisture greatly influences vegetation response, root growth, excavation difficulty, albedo, trafficability, construction, conductivity, soil chemical interactions, workability, chemical transport, strength, shrinking and swelling, frost action, seed germination, and many other properties, qualities, and interpretations. Soil moisture states are significant to soil taxonomic classification, wetland classification, and other classification systems. The recording of soil moisture states helps to document the soil classification as well as convey information useful for crop and land management models. Measurement Soil water status can be measured using tensiometers or moisture tension plates. Soil water status also can be field estimated. Chapter 3 of the Soil Survey Manual provides more information. It is important to note that the three water state classes and eight subclasses described in the Soil Survey Manual are used to describe the moisture state at a point in time for individual pedons (spatial and temporal point data), while the water state classes discussed here are used to estimate the mean monthly aggregated moisture conditions for a map unit component. As a consequence, only three classes are used and the definitions for the moist and wet classes are modified from the definitions in the Soil Survey Manual. The wet class used here includes only the satiated wet class and corresponds to a free water table. The moist class is expanded to include the nonsatiated wet class given in the Soil Survey Manual. Dry is separated from moist at 15 bar suction. Wet satiated has a tension of 0.0 bar or less (zero or positive pore pressure). Changes in natural patterns of water movement from dams and levees are considered in evaluating and entering soil moisture status. Infiltration, saturated hydraulic conductivity, and organic matter, which affect soil moisture movement, are strongly impacted by land cover and land use. Land use and land cover should be considered as a mapping tool for separating map units or map unit components. The difference in soil moisture status resulting from differences in land use and land cover constitute a difference in soil properties. However, conservation practices, such as irrigating and fallowing the land, alter the soil moisture status but are not considered in the map unit component data. Use-dependent databases may allow entries for these altered states in the future. Permanent installations, such as drainage ditches and tile, affect soil moisture status, and the drained condition should be reflected in the soil moisture status entries for map unit components that are mapped as “drained.” Undrained areas are mapped as “undrained” components, and the entries for soil moisture status reflect the undrained condition. Irrigation and drainage canals are shown on soil maps; their effects on the soil should be shown in the properties of the soils in mapping and in the property records. Soils that are now wet because of excessive irrigation and leaking canals should be mapped, and their properties should reflect the current soil moisture status. Guiding Concepts The intent is to describe a mean moisture condition, by month, for a soil component. Layer depths may or may not be the same as horizon depths in the Component Horizon Table. Layers define the zone having a specific soil moisture state. If the soil is wet throughout 0 to 200 cm, then one entry (“wet”) is made for 0 to 200 cm for that month. For frozen soils, enter the appropriate soil moisture state that the soil would have if thawed. For example, if the soil is frozen and then determined to be wet when thawed, enter “wet.” The horizons can be subdivided or combined, as appropriate, into layers for the various soil moisture states as needed. Remember that these are monthly averages for the extent of the component across the landscape. The entries are expected to come from the best estimates that local knowledge can provide. If local knowledge is supported by data, so much the better. The information as aggregated data is not expected to be exact but should be generalized and reflect an average condition. Entries for the representative values (rv) on distance to the upper and lower boundary of the moisture layer should reflect the soil moisture conditions expected in a normal year, as defined in the latest edition of the Keys to Soil Taxonomy. Make entries for each month by layer. Enter the dominant condition for the month. This is the condition that exists for more than 15 days on the long-term average. The low and high values represent the depth range within the component for the normal year; they should not represent the extremes, such as years of drought. If the depth to free water fluctuates during the month, use the depth for the average between the high and low levels. Entries.—Enter the soil moisture status as dry, moist, or wet for each soil layer for each month. Enter only one soil moisture state for a given layer. The number of layers depends upon the number of changes of soil moisture status in the profile. Enter the values for component soil moisture depth to top and depth to bottom that represents the distance, in whole centimeters, from the soil surface to the top and bottom respectively, of each soil layer for each month. Part 618, Subpart B, Exhibits, Section 618.97 contains examples of entries in a worksheet format that graphs soil moisture status by month and depth. 618.61 Soil Slippage Potential Definition.—“Soil slippage potential” is the hazard that a mass of soil will slip when vegetation is removed, soil water is at or near saturation, and other normal practices are applied. Conditions that increase the hazard of slippage but are not considered in this rating are undercutting lower portions or loading the upper parts of a slope or altering the drainage or offsite water contribution to the site, such as through irrigation. The publication Landslides Investigation and Mitigation Special Report 247 (Transportation Research Board, National Research Council, 1996) provides additional information on landscape slippage. Significance.—Slippage is an important consideration for engineering practices, such as constructing roads and buildings, and for forestry practices. Estimates.— Soil slippage potential classes are estimated by observing slope; lithology, including contrasting lithologies; strike and dip; surface drainage patterns; and occurrences of such features as slip scars and slumps. Entries.—Enter one of the following soil slippage potential classes for the component: High (unstable) Medium (moderately unstable) Low (slightly unstable to stable) 618.62 Soil Temperature Definition.—“Soil temperature” is the temperature calculated as both the mean annual temperature at a single depth in the soil and the mean monthly temperature calculated at a specified depth range for each month of the year. Estimates.— Soil temperature according to depth can be estimated from measured soil temperatures of the vicinity. Air temperature fluctuations, soil moisture, aspect, slope, color, snow cover, plant cover, and residue cover affect soil temperature. Estimates of soil temperature should take these factors into account when soil temperatures are extrapolated from one soil map unit component to another. Measurement.— Soil temperature can be measured by many types of thermometers, including mercury, bimetallic, thermisters, and thermocouples. Many types of thermometers can be configured for remote, unattended operation. Mean Annual Soil Temperature (MAST) Definition.—“Mean annual soil temperature (MAST)” is the temperature generally determined at a depth of 50 cm below the soil surface, or at the upper boundary of a root-limiting layer as defined in Soil Taxonomy, whichever is shallower. Entries.—Enter the high, low, and representative values for the range of mean annual soil temperature for the component as the long-term average of the mean monthly soil temperatures in the Component table. The long-term average is generally considered to be a 30-year average. The range of valid entries is from -40 to 50 degrees Celsius, and tenths (one decimal place) are allowed. Mean Monthly Soil Temperature Definition.—“Mean monthly soil temperature” is the long-term monthly average of the mean daily high and daily low soil temperature at a specified depth for the month in question. Long-term is generally considered to be a 30-year average. Entries.—Enter soil temperature for the component as the long-term monthly average of the mean daily soil temperature at a specified depth for the month in question in the Component Soil Temperature table. The long-term average is generally considered to be a 30-year average. The range of valid entries is from -25 to 50 degrees Celsius, and only whole numbers (integers) are allowed. The number of layers populated depends upon the number of changes of soil temperature status in the profile. Soil Temperature, Depth to Top Definition.—“Soil temperature, depth to top” is the distance from the top of the soil to the upper boundary of the soil temperature layer. Entries.—Enter the value for soil temperature depth to top that represents the distance, in centimeters, from the soil surface to the top of each soil temperature layer for each month in the Component Soil Temperature table. Soil Temperature, Depth to Bottom Definition.—“Soil temperature, depth to bottom” is the distance from the top of the soil to the lower boundary of the soil temperature layer. Entries.—Enter the value for soil temperature depth to bottom that represents the distance, in centimeters, from the soil surface to the bottom of each soil temperature layer for each month in the Component Soil Temperature table. 618.63 Subsidence, Initial and Total Definition.—“Subsidence” is the decrease in surface elevation as a result of the drainage of wet soils that have organic layers or semifluid, mineral layers. Initial subsidence is the decrease of surface elevation that occurs within the first 3 years of the drainage of these wet soils. Total subsidence is the potential decrease of surface elevation as a result of the drainage of these wet soils. Significance The susceptibility of soils to subsidence is an important consideration for organic soils that are drained. If these soils are drained for community development, special foundations are needed for buildings. Utility lines, sidewalks, and roads that lack special foundations may settle at different rates, thus causing breakage, high maintenance costs, and inconvenience. If the soils are drained for farming, the long-term effects of subsidence, the possible destruction of land if it subsides below the water table, and possible legal implications where the soils are in wetlands must be considered. Subsidence as a result of drainage is attributed to the following factors. The first three factors are responsible for the initial subsidence that occurs rapidly, specifically within about 3 years after the water table is lowered. Shrinkage from drying Consolidation because of the loss of ground-water buoyancy. Compaction from tillage or manipulation Wind erosion Burning Biochemical oxidation After the initial subsidence, a degree of stability is reached and the loss of elevation declines to a steady rate, primarily because of oxidation. The oxidation and subsidence continue at this slower rate until stopped by the water table or underlying mineral material. The rate of subsidence depends on— Ground-water depth. Amount of organic matter. Kind of organic matter. Soil temperature. pH. Biochemical activity. Estimates A number of studies have been made to measure actual subsidence. Other useful studies have measured the bulk density of organic soils after drainage. Based on these studies, some general guidelines can be given for initial and total subsidence. Initial subsidence generally is about half of the depth to the lowered water table or to mineral soil, whichever is shallower. It occurs within about 3 years after drainage. Total subsidence is the total depth to the water table or the thickness of the organic layer, whichever is shallower. It is rarely reached, except where organic layers are thin or where drainage systems have been installed for a long time. Measurement.— After organic soils have been drained and cultivated for a number of years, they reach a nearly steady rate of subsidence that is reflected by the rather stable bulk density. Unpublished studies by the NSSC Kellogg Soil Survey Laboratory have shown that the bulk density of the organic component, such as that with the percent mineral calculated out, stabilizes at around 0.27 g/cc for surface layers and 0.18 g/cc for subsurface layers. These values can be used to calculate the amount of subsidence at some time in the future as compared to the thickness of soil at the time of observation or measurement. The procedure is as follows: Sample the surface and subsurface layers for field state bulk density. Methods are described in the Handbook of Soil Survey Investigations Field Procedures, I 4-2, 1971, USDA, Soil Conservation Service, and in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Calculate out the weight contribution of the mineral component to obtain the bulk density of the organic component (DbOM). This manipulation allows bulk densities to be on a common base so that various layers can be compared. The formula for the computation is as follows: DbOM = Db (1 - percent mineral/100), where Db is the field state bulk density. Calculate the subsidence percent (SP) for surface and subsoil horizons as follows: For surface horizons: SP = 100 - [(DbOM/0.27) x 100] For subsurface horizons: SP = 100 - [(DbOM/0.18) x 100] Where DbOM is obtained from step (2). Convert initial subsidence percent to depth of subsidence in inches as follows: Entries.—Enter the high, low, and representative values that represent the range for initial and total subsidence, in centimeters, for the map unit component. The range of valid entries is from 0 to 999, and only whole numbers (integers) are allowed. If subsidence is not a concern, enter “0.” 618.64 Sum of Bases Definition.—“Sum of bases” is the sum of the basic cations calcium, magnesium, potassium, and sodium that are extractable from the < 2 mm soil fraction using a solution of ammonium acetate (NH4OAc, pH 7). Significance.—Sum of bases is important for certain evaluations of soil nutrient availability or of the effect of waste additions to the soil. Sum of extractable bases is used directly in soil classification as a criterion to classify soils in most of the Eutric subgroups of Andisols. It is also used indirectly in soil classification to calculate percent base saturation by the sum of cations method. Base saturation by sum of cations is used as a criterion for Ultisols, Ultic subgroups of Alfisols, Andisols, and Mollisols, Alfic and Dystric subgroups of Inceptisols, and Alfic subgroups of Spodosols. Measurement.—Sum of bases is calculated from the results of methods outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Sum of bases is reported in centimoles per kilogram (cmol(+) kg-1), which are equivalent to milliequivalents per 100 grams (meq 100 g-1) of soil. Entries.—Enter the range of sum of bases as milliequivalents per 100 grams (meq 100 g-1) of soil for the horizon. The range of valid entries is from 0 to 300 and tenths (one decimal place) are allowed. 618.65 Surface Fragments Definition.—“Surface fragments” are unattached, cemented pieces of bedrock, bedrocklike material, durinodes, concretions, nodules, or pedogenic horizons (e.g., petrocalcic fragments) 2 mm or larger in diameter and woody material 20 mm or larger in diameter that are exposed at the surface of the soil. Surface fragments can be rock fragments, pararock fragments, or wood fragments, as defined in Section 618.31. Vegetal material other than wood fragments, whether live or dead, is not included. Surface Fragment Cover Percent Definition.—“Surface fragment cover percent” is the percent of ground covered by fragments 2 mm or larger in diameter (20 mm or larger in diameter for wood fragments). Significance.—Fragments on the soil surface are used as map unit phase criteria and greatly affect the use and management of the soil. They affect equipment use, erosion, excavation, and construction. They act as a mulch, slowing evaporation and armoring the soil against rainfall impact. They also affect the heating and cooling of soils. Estimates.— An estimation of cover by surface fragments can be made visually without quantitative measurement, by transect techniques, or by some combination of visual and quantitative measures. Chapter 3 of the Soil Survey Manual provides more information. Entries.—Enter the high, low, and representative values for the percent of the surface covered by each size class and kind of fragment populated in the Component Surface Fragments Table of the NASIS database. The range of valid entries is from 0 to 100 percent, and hundredths (two decimal places) are allowed. Surface Fragment Kind Definition.—“Surface fragment kind” is the lithology or composition of the surface fragments 2mm or larger in diameter (20 mm or larger in diameter for wood fragments). Significance.—Fragments vary according to their resistance to weathering. Consequently, fragments of some lithologies are more suited than others for use as building stone, road building material, or riprap to face dams and stream channels. Entries.—Enter the appropriate fragment kind name for the record of fragments populated in the Component Surface Fragments Table in the NASIS database. The class names are present in a choice list and can also be viewed in the NASIS data dictionary. Surface Fragment Size Definition.—“Surface fragment size” is the size based on the multiaxial dimensions of the surface fragments. Significance.—The size of surface fragments is significant to the use and management of the soil. The adjective form of fragment size is used as phase criteria for naming map units. The size affects equipment use, excavation, construction, and recreational uses. Classes Classes of surface fragment size are subdivided based on the shape of the fragments (described below). Flat fragment classes are: Flat fragment class Length of fragment (mm) Channers 2-150 Flagstones 150-380 Stones 380-600 Boulders ≥600 Nonflat fragment classes are: Nonflat fragment class Diameter (mm) Gravel 2-75 Fine gravel 2-5 Medium gravel 5-20 Coarse gravel 20-75 Cobbles 75-250 Stones 250-600 Boulders ≥600 Gravel is a collection of fragments having a diameter ranging from 2 to 75 mm. Individual fragments in this size range are properly referred to as pebbles, not “gravels.” For fragments that are less than strongly cemented, “para” is used as a prefix to the above terms, i.e., paracobbles. Entries.—Enter the high, low, and representative values for each size class populated in the Component Surface Fragments Table in the NASIS database. Valid entries are 2 millimeters (mm) or larger, and only whole numbers (integers) are allowed. Mean Distance Between Rocks Definition.—“Mean distance between rocks” is the average distance between surface stones, boulders, or both, measured between edges. Significance.—The mean distance between rocks is a field clue for naming stony or bouldery map units. The closer the distance, the more equipment limitations there are for harvesting forestland or soil cultivation. Estimates.— Table 3-12 of the Soil Survey Manual shows the distance between stones and boulders if the diameter is 0.25 m, 0.6 m, or 1.2 m. This table should be used with caution because stones and boulders are rarely equally spaced or have the same diameter. Entries.—Enter the high, low, and representative values for the mean distance between rocks. The range of valid entries is from 0 to 50 meters, and hundredths (two decimal places) are allowed. Surface Fragment Roundness Definition.—“Surface fragment roundness” is an expression of the sharpness of edges and corners of surface fragments. Classes.—The surface fragment roundness classes are: Roundness class Definition Very angular Strongly developed faces with very sharp, broken edges Angular Strongly developed faces with sharp edges (SSM) Subangular Detectable flat faces with slightly rounded corners Subrounded Detectable flat faces with well-rounded corners (SSM) Rounded Flat faces absent or nearly absent with all corners rounded (SSM) Well rounded Flat faces absent with all corners rounded Entries.—Enter the appropriate surface fragment roundness class name for the record of surface fragments populated in the Component Surface Fragments Table in the NASIS database. Surface Fragment Hardness Definition.—“Surface fragment hardness” is equivalent to the rupture resistance cemented of a surface fragment of specified size that has been air-dried and then submerged in water. Measurements.—Procedures and classes of cementation are listed with the rupture resistance classes in the Soil Survey Manual. Classes are described for similar specimens about 25-30 mm on edge which are air-dried and then submerged in water for at least 1 hour. The specimen is compressed between extended thumb and forefinger, between both hands, or between the foot and a hard flat surface. If the specimen resists compression, a weight is dropped onto it from progressively greater heights until it ruptures. Failure is considered at the initial detection of deformation or rupture. Stress applied in the hand should be over a 1-second period. The tactile sense of the class limits may be learned by applying force to top-loading scales and sensing the pressure through the tips of the fingers or through the ball of the foot. Postal scales may be used for the resistance range that is testable with the fingers. A bathroom scale may be used for the higher rupture resistance range. Significance.—The hardness of a surface fragment is significant where the rupture resistance class is strongly cemented or greater. These classes can impede or restrict the movement of soil water vertically through the soil profile and have a direct impact on the quality and quantity of ground water and surface water. Entries.—Enter the appropriate class name for each record of surface fragments populated in the Component Surface Fragments Table in the NASIS database. Choose the term without the word “cemented” (e.g., choose the “moderately” class to represent the moderately cemented class). Surface Fragment Shape Definition.—“Surface fragment shape” is a description of the overall shape of the surface fragment. Classes.—The surface fragment shape classes are “flat” and “nonflat.” Entries.—Enter the appropriate surface fragment shape class name for each record of surface fragments populated in the Component Surface Fragments Table in the NASIS database. 618.66 T Factor Definition.—The “T factor” is the soil loss tolerance (in tons per acre). It is defined as the maximum amount of erosion at which the quality of a soil as a medium for plant growth can be maintained. This quality of the soil to be maintained is threefold in focus. It includes maintaining the surface soil as a seedbed for plants, the atmosphere-soil interface to allow the entry of air and water into the soil and still protect the underlying soil from wind and water erosion, and the total soil volume as a reservoir for water and plant nutrients, which is preserved by minimizing soil loss. Erosion losses are estimated by USLE and RUSLE2. Classes.—The classes of T factors are 1, 2, 3, 4, and 5. Significance.—Soil loss tolerances commonly serve as objectives for conservation planning on farms. These objectives assist in the identification of cropping sequences and management systems that can maximize production and also sustain long-term productivity. T factors represent the goal for maximum annual soil loss. Guidelines.—Conservation objectives for soil loss tolerance include maintaining a suitable seedbed and nutrient supply in the surface soil, maintaining an adequate depth and quality of the rooting zone, and minimizing unfavorable changes in water status throughout the soil. A single T factor is assigned to each map unit component. Estimates.— The T factor is assigned to soils without respect to land use or cover. T factors are assigned to compare soils and do not imply differences to vegetation response directly. Many of the factors used to assign a T factor are also important to vegetation response, but the T factor is not assigned to imply vegetation sensitivity to all vegetation. The general guideline given in Part 618, Subpart B, Exhibits, Section 618.93 is used to assign T factors but more specific criteria are used to select limiting soil properties. Entries.—The estimated soil loss tolerance should be calculated from the soil properties and qualities posted in the database for each map unit component based generally on the guideline given in Part 618, Subpart B, Exhibits, Section 618.93. Acceptable values are 1, 2, 3, 4, and 5. 618.67 Taxonomic Family Temperature Class Definition The soil temperature classes are part of the family categorical level as defined in Soil Taxonomy. They differ from “soil temperature regimes,” (Data Element: taxonomic temp regime), in that the cryic temperature regime is divided between the frigid and isofrigid classes based on differences in mean winter and mean summer soil temperatures. Soil temperature classes are based on mean annual and mean seasonal soil temperatures using the Celsius (centigrade) scale and taken either at a depth of 50 cm from the soil surface or at a lithic or paralithic contact, whichever is shallower. For soil families in Gelisols, Gelic suborders, and Gelic great groups the soil temperature classes, defined in terms of the mean annual soil temperature, are as follows: Hypergelic: -10°C or lower Pergelic: -4°C to -10°C Subgelic: +1°C to -4°C For soil families that have a difference of 6°C or more between mean summer (June, July, and August in the northern hemisphere) temperature and mean winter (December, January, and February in the northern hemisphere) temperature, the soil temperature classes, defined in terms of the mean annual soil temperature, are as follows: Frigid: Lower than 8°C Mesic: 8°C to 15°C Thermic: 15°C to 22°C Hyperthermic: 22°C or higher For soil families that have a difference of less than 6°C between the mean summer and mean winter soil temperatures, the soil temperature classes, defined in terms of the mean annual soil temperature, are as follows: Isofrigid: Lower than 8°C Isomesic: 8°C to 15°C Isothermic: 15°C to 22°C Isohyperthermic: 22°C or higher Significance.—All soils have a taxonomic soil temperature class. Soil temperature classes are used as family differentiae in all the orders defined in Soil Taxonomy. The names are used as part of the family name unless the criteria for a higher taxon carry the same limitation. The frigid or isofrigid class is implied in all cryic suborders and great groups, but the class is not used as part of the family name because it would be redundant. Estimates.— Estimates of soil temperature classes are made with models that use climatic data including mean annual and mean seasonal air temperatures, precipitation, and evapotranspiration. Some models include snow cover, topographic, and vegetative inputs. Measurement.— The Celsius (centigrade) scale is the standard. It is assumed that the temperature is that of a nonirrigated soil. The soil temperature classes are based on long-term averages of mean annual and mean seasonal soil temperatures taken either at a depth of 50 cm from the soil surface or at a lithic or paralithic contact, whichever is shallower. Entries.—Enter the appropriate soil temperature class from the following list: Frigid Hypergelic Hyperthermic Isofrigid Isohyperthermic Isomesic Isothermic Mesic Pergelic Subgelic Thermic Not used 618.68 Taxonomic Moisture Class Definition.—Soil moisture classes refer to the soil moisture regimes defined in Soil Taxonomy. Soil moisture regimes are defined by the presence or absence either of ground water or of water held at a tension of less than 1500 kPa, in the soil or in specific horizons, by periods of the year. Significance.—All soils have a soil moisture regime. Soil moisture regimes are used as differentiae in all the orders defined in Soil Taxonomy. Data on the moisture regime are used for making interpretations for cropland agriculture, correlating soils to ecological sites, and determining suitability for wildlife habitat. The moisture regime of some soils is not apparent in the classification given in Soil Taxonomy. Ustolls and Xerolls, for example, can have an aridic moisture regime. Some soils have more than one moisture regime. An example is a soil that meets the requirements of the aquic moisture regime in the wet season and also meets the requirements of the ustic regime. Estimates.— Estimates of soil moisture regimes are made with models that use climatic data, including mean annual and mean seasonal air temperatures, precipitation, and evapotranspiration. Some models include topographic and vegetative inputs. The soil moisture control section, also defined in Soil Taxonomy, is used to facilitate the estimation of soil moisture regimes. See Soil Survey Technical Note 9 for more guidance. Measurement.— The soil moisture regimes are based on annual and seasonal soil moisture measurements taken in the soil moisture control section. The soil should not be irrigated, fallowed, or influenced by other moisture-altering practices. Entries.—Enter the appropriate soil moisture regimes from the following list: Definition.—“Subclasses of soil moisture regimes” are defined at the subgroup categorical level in Soil Taxonomy. The criteria differ among the great groups. For example aquic, aridic, and udic are subclasses of the soil moisture regime in Haplustalfs. A subclass is entered for all soils in a great group that meet the subclass criteria, even if the subclass is not part of the taxonomic classification. For example, aquic, aridic, udic, or typic should be used as a subclass of the soil moisture regime in Lithic Haplustalfs if the criteria are met. Significance.—Subclasses of soil moisture regimes are used at the subgroup categorical level in all orders in Soil Taxonomy except Histosols. They typically indicate an intergrade between two moisture regimes that affect the use and management of the soil. The subclasses of soil moisture regimes are used for making interpretations for cropland agriculture, correlating soils to ecological sites, and determining suitability for wildlife habitat. Estimates.— Estimates of subclasses of soil moisture regimes are made with models that use climatic data, including mean annual and mean seasonal air temperatures, precipitation, and evapotranspiration. Some models include topographic and vegetative inputs. The soil moisture control section, also defined in Soil Taxonomy, is used to facilitate estimation of some subclasses of soil moisture regimes. For more guidance, see Soil Survey Technical Note 9. Measurement.— The subclasses of soil moisture regimes are based on annual and seasonal soil moisture measurements taken in the soil moisture control section. The soil should not be irrigated, fallowed, or influenced by other moisture-altering practices. Entries.—Enter the appropriate subclass of soil moisture regimes from the following list: Aeric Anthraquic Aquic Aridic (torric) Oxyaquic Typic Udic Ustic Xeric 618.70 Taxonomic Temperature Regime (Soil Temperature Regimes) Definition.—“Soil temperature regimes” refer to the temperature regimes as defined in Soil Taxonomy. Significance.—Soil temperature regimes are used as differentiae above the family categorical level in all orders in Soil Taxonomy. (Soil temperature classes, defined above, are used as family differentiae.) Soil temperature regimes greatly affect the use and management of soils, particularly the selection of adapted plants. Temperature regimes are used for making interpretations for cropland agriculture, correlating soils to ecological sites, and determining suitability for wildlife habitat. Estimates.— Estimates of soil temperature regimes are made with models that use climatic data including mean annual and mean seasonal air temperatures, precipitation, and evapotranspiration. Some models include topographic and vegetative inputs. Measurement.— The soil temperature regime is based on mean annual and seasonal soil temperatures using the Celsius (centigrade) scale and taken either at a depth of 50 cm from the soil surface or at a lithic or paralithic contact, whichever is shallower. Entries.—Enter the appropriate soil temperature regimes from the following list: Gelic Cryic Frigid Mesic Thermic Hyperthermic Isofrigid Isomesic Isothermic Isohyperthermic 618.71 Texture Class, Texture Modifier, and Terms Used in Lieu of Texture Definition.—“Texture class” refers to the soil texture classification used by the U.S. Department of Agriculture as defined in the Soil Survey Manual. Soil texture is the relative proportion, by weight, of the particle separate classes finer than 2 mm in equivalent diameter. The material finer than 2 mm is the fine-earth fraction. Material 2 mm or larger is rock or pararock fragments. Significance.—Soil texture influences engineering works and plant growth and indicates how soils formed. Soil texture has a strong influence on soil mechanics and the behavior of soil when it is used as construction or foundation material. It influences such engineering properties as bearing strength, compressibility, saturated hydraulic conductivity, shrink-swell potential, and compaction. Engineers are also particularly interested in rock and pararock fragments. Soil texture influences plant growth by its affect on aeration, the water intake rate, the available water capacity, the cation-exchange capacity, saturated hydraulic conductivity, erodibility, and workability. Changes in texture as related to depth are indicators of how soils formed. When texture is plotted with depth, smooth curves indicate translocation and accumulation. Irregular changes in particle-size distribution, especially in the sand fraction, may indicate lithologic discontinuities, specifically differences in parent material. Soil texture is used for soil classification in criteria for certain taxa such as the Psamments suborder, “Psamm” great groups, and Arenic, Grossarenic, and Psammentic subgroups. Soil texture is also used in the family category of Soil Taxonomy for differentiae such as particle-size class. Measurement.— USDA texture can be measured in the laboratory by determining the proportion of the various size particles in a soil sample. The analytical procedure is called particle-size analysis or mechanical analysis. Stone, gravel, and other material 2 mm or larger are sieved out of the sample and thus are not considered in the analysis of the sample. Their amounts are measured separately. Of the remaining material smaller than 2 mm, the amount of the various sizes of sand is determined by sieving. The amount of silt and clay is determined by a differential rate of settling in water. Either the pipette or hydrometer method is used for the silt and clay analysis. Organic matter and dissolved mineral matter are removed in the pipette procedure but not in the hydrometer procedure. The two procedures are generally very similar, but a few samples, especially those with high organic matter or high soluble salts, exhibit wide discrepancies. The detailed procedures are outlined in Soil Survey Investigations Report No. 42, Soil Survey Laboratory Methods Manual, Version 4.0, November 2004, USDA, NRCS. Estimates The determination of soil texture for the less than 2 mm material is made in the field mainly by feeling the soil with the fingers. The soil must be well moistened and rubbed vigorously between the fingers for a proper determination of texture class by feel. This method requires skill and experience but good accuracy can be obtained if the field soil scientist frequently checks his or her estimates against laboratory results. Many NRCS offices collect reference samples for this purpose. The content of particles larger than 2 mm cannot be evaluated by feel. The content of the fragments is determined by estimating the proportion of the soil volume that they occupy. Fragments in the soil are discussed in Section 618.31. Each soil scientist must develop the ability to determine soil texture by feel for each genetic soil group according to the standards established by particle-size analysis. Soil scientists must remember that soil horizons that are in the same texture class but are in different subgroups or families may have a different feel. For example, natric horizons generally feel higher in clay than “non-natric” horizons. Laboratory analysis generally shows that the clay in natric horizons is less than the amount estimated from the field method. The scientist needs to adjust judgment and not the size distribution standards. Entries.—Texture is displayed by the use of six data elements in the NASIS database: texture class, texture modifier, texture modifier and class, stratified texture flag, representative value indicator, and terms used in lieu of texture. Only use multiple textures if they interpret the same for the horizon. Only textures that represent complete horizons should be entered. In NASIS the representative value indicator is identified (i.e., representative? = yes) for the single row that contains the texture term considered typical for each interpretive horizon of the component. This choice should match the representative values of the various soil particle-size separates posted elsewhere in the database. Texture Class Definition “Texture class” is an expression, based on the USDA system of particle sizes, for the relative portions of the various size groups of individual mineral soil grains less than 2 mm equivalent diameter in a mass of soil. Each texture class has defined limits for each particle separate class of mineral particles less than 2 mm in effective diameter. The basic texture classes, in the approximate order of increasing proportions of fine particles, are sand, loamy sand, sandy loam, loam, silt loam, silt, sandy clay loam, clay loam, silty clay loam, sandy clay, silty clay, and clay. The sand, loamy sand, and sandy loam classes may be further subdivided into coarse, fine, or very fine. The basic USDA texture classes are given graphically in Part 618, Subpart B, Exhibits, Section 618.87 as a percentage of sand, silt, and clay. The chart at the bottom of the figure shows the relationship between the particle size and texture classes among the AASHTO, USDA, and Unified soil classification systems. Definition.—“Terms used in lieu of texture” are substitute terms applied to materials that do not fit into a texture class because of high organic matter content, high fragment content, high gypsum content, cementation, or another reason. Examples include artifacts, bedrock, gravel, and muck. Part 618, Subpart B, Exhibits, Section 618.94 provides a list of these terms and their codes. Some of these terms may be modified with terms from the list of texture modifiers, such as mossy (code MS) when used to modify the term peat (i.e., “mossy peat”). Application The terms used in lieu of texture “highly decomposed plant material,” “moderately decomposed plant material,” and “slightly decomposed plant material” (codes HPM, MPM, and SPM), should only be used to describe near surface horizons composed of plant material in various stages of decomposition that are saturated with water for less than 30 cumulative days in normal years and are not artificially drained. These terms are used to describe folistic epipedons (i.e., in mineral soils only) or organic horizons of any thickness (i.e., in organic or mineral soils) provided they meet the saturation requirements. The terms “muck,” “mucky peat,” and “peat” (codes MUCK, MPT, and PEAT) are used to describe histic epipedons (i.e., in mineral soils only) and organic horizons of any thickness (i.e., in organic or mineral soils) that are saturated with water for 30 or more cumulative days in normal years or are artificially drained. For soil materials with 40 percent or more, by weight, gypsum in the fine-earth fraction, gypsum dominates the physical and chemical properties of the soil and particle-size classes are not meaningful. Two terms in lieu of texture are used. “Coarse gypsum material” (code CGM) is used for these materials where 50 percent or more of the fine-earth fraction is comprised of particles ranging from 0.1 to 2.0 mm in diameter. “Fine gypsum material” (code FGM) is applied to materials where less than 50 percent of the fine-earth fraction is comprised of particles ranging from 0.1 to 2.0 mm in diameter. The term “material,” (code MAT), is generic and requires the use of a texture modifier. It is intended for cemented diagnostic horizons such as duripans, petrocalcic horizons, and petrogypsic horizons (coded CEM-MAT), by using the texture modifier “cemented” with the term in lieu of texture “material.” The concatenated texture term for such horizons in pedon descriptions is “cemented material.” In the past, texture modifier terms, such as “coprogenous,” “gypsiferous,” and “marly,” were used to describe material, but such use has been discontinued and is no longer permitted. Examples of current usage are shown below and combine the texture modifier with an appropriate texture class (e.g., marly silt loam). Definition.—“Texture modifier” is a term used to denote the presence of a condition or object other than sand, silt, or clay. Application.—Texture modifier terms may apply to both texture and terms used in lieu of texture. Some may apply to both, others only apply to one or the other. Combinations of some texture modifiers are allowed. A list of allowable texture modifier terms and their codes is given in Part 618, Subpart B, Exhibits, Section 618.94. Some rules of application are given below. If the content of fragments equals 15 percent or more, by volume, texture modifiers are used. An example is gravelly loam or parachannery loam. The adjectives “very” and “extremely” are used when the content of fragments equals 35 to less than 60 percent and 60 to less than 90 percent, by volume, respectively. Texture modifiers, such as paragravelly and paracobbly, are used to identify the presence of pararock fragments. The size, shape, and amounts of pararock fragments required for these terms are the same as for rock fragments. “Mucky” and “peaty” are used to modify near surface horizons of mineral soils that are saturated with water for 30 or more cumulative days in normal years or are artificially drained. An example is mucky loam. Excluding live roots, the horizon has an organic carbon content (by weight) of one of the following: 5 to < 12 percent if the mineral fraction contains no clay 12 to < 18 percent if the mineral fraction contains 60 percent or more clay “Highly organic” is used to modify near surface horizons of mineral soils that are saturated with water for less than 30 cumulative days in normal years and are not artificially drained. Excluding live roots, the horizon has an organic carbon content (by weight) of one of the following: 5 to < 20 percent if the mineral fraction contains no clay 12 to < 20 percent if the mineral fraction contains 60 percent or more clay Compound texture modifiers may be used. For example, a term may be used to indicate the presence of fragments and another used to indicate some nonfragment condition. The term used to indicate fragments should be listed first. Examples are very gravelly mucky silt loam and paragravelly ashy loam. In some instances, mineral soil may contain a combination of both artifacts and fragments in the soil such as rock fragments and pararock fragments. In all cases, the artifacts, rock fragments, and pararock fragments are each described separately. The assignment of texture modifiers for such horizons is handled differently depending on the nature of the artifacts. Artifacts in soils which are discrete (i.e., ≥ 2 mm), cohesive, and persistent (e.g., concrete) function in a manner which is similar to rock fragments. Artifacts which are either noncohesive or nonpersistent (e.g., cardboard) behave differently than other discrete artifacts and also rock fragments. When describing the texture of soil horizons with artifacts or a combination of artifacts and fragments, the following rules of application are followed: Describe the individual kinds and amounts (percent by volume) of artifacts and any fragments, if present. Record all pertinent attributes for artifacts (see Section 618.5), paying particular attention to data on artifact cohesion and persistence. If the combined volume of artifacts, which are both cohesive and persistent, plus any rock fragments present is less than 15 percent, use the following table: Less than 15 percent: No artifact texture modifier is used. 15 to < 35 percent: The adjectival term “artifactual” is used as a modifier of the texture class, such as “artifactual loam.” 35 to < 60 percent: The adjectival term “very artifactual” is used as a modifier of the texture class, such as “very artifactual loam.” 60 percent to < 90 percent: The adjectival term “extremely artifactual” is used as a modifier of the texture class, such as “extremely artifactual loam.” 90 percent or more No texture modifier terms are used. If there is too little fine earth to determine the texture class (less than about 10 percent, by volume) the term used in lieu of texture, “artifacts”, is populated. If both artifacts and rock fragments are present and the combined volume of rock fragments and artifacts, which are both cohesive and persistent, is 15 percent or more, assign dual rock fragment-artifact texture modifiers. Dual rock fragment-artifact texture modifiers are based on the combined volume of both. The modifiers are concatenated terms joined with a hyphen. For example, use “gravelly-artifactual loam” as the texture modifier for a horizon with a fine-earth texture class of loam that contains 10 percent quartzite gravel, 3 percent brick (a cohesive and persistent artifact), 2 percent glass (a cohesive and persistent artifact), and 25 percent plasterboard (a noncohesive artifact). See Subpart B, Exhibits, Section 618.94 for the list of 18 dual rock fragment-artifact texture modifiers. If artifacts, pararock fragments, and rock fragments are present, but the combined volume of artifacts, which are both cohesive and persistent, and any rock fragments present is less than 15 percent, compound texture modifiers are used. The compound texture modifiers connote only the artifacts and the pararock fragments. The modifier for artifacts is assigned (using the table shown above) preceding the texture modifier for pararock fragments. Some examples are “artifactual paracobbly coarse sandy loam” for a horizon that contains 20 percent rubber (e.g., shredded tires) and 20 percent granite paracobbles and “very artifactual parachannery clay” for a horizon with 40 percent carpet pieces and 20 percent siltstone parachanners. If a horizon includes both rock fragments and pararock fragments, use the following rules for selecting texture modifiers: Describe the individual kinds and amounts of rock fragments and pararock fragments. Do not use a fragment texture modifier when the combined volume of rock fragments and pararock fragments is less than 15 percent. When the combined volume of rock fragments and pararock fragments is 15 percent or more and the volume of rock fragments is less than 15 percent, assign pararock fragment modifiers based on the combined volume of fragments. For example, use paragravelly as a texture modifier for soils with 10 percent rock and 10 percent pararock gravel-sized fragments. When the volume of rock fragments is 15 percent or more, use the appropriate texture modifier for rock fragments (see Part 618, Subpart B, Exhibits, Section 618.90), regardless of the volume of pararock fragments. (Do not add the volume of rock and pararock fragments to determine the texture modifier.) The definitions of the following four compositional texture modifiers guide their usage: Hydrous.—Material that has andic soil properties and an undried 15 bar (1500 kPa)water content of 100 percent or more of the dry weight. Medial.—Material that has andic soil properties and has a 15 bar (1500 kPa) water content of less than 100 percent on undried samples and of 12 percent or more on air-dried samples. Ashy.—Material that has andic soil properties and is neither hydrous nor medial or material that does not have andic soil properties and the fine-earth fraction contains 30 percent or more particles 0.02 to 2.0 mm in diameter, of which 5 percent or more is composed of volcanic glass and the [(aluminum plus 1/2 iron percent by ammonium oxalate) times 60] plus the volcanic glass percent is equal to or more than 30. Woody, grassy, mossy, and herbaceous texture modifiers are only used to modify muck, peat, or mucky peat terms (used for histic epipedons and organic horizons of any thickness that are saturated with water for 30 or more cumulative days in normal years, or are artificially drained, including those in Histels and Histosols, except for Folists). The definitions of the following four compositional texture modifiers guide their usage: Woody.—Any material that contains 15 percent or more wood fragments larger than 2 cm in size or organic soil materials other than SPM, MPM, or HPM, that contain 15 percent or more fibers that can be identified as wood origin and contain more wood fibers than any other kind of fiber. Grassy.—Organic soil material that contains more than 15 percent fibers that can be identified as grass, sedges, cattails, and other grasslike plants and contains more grassy fibers than any other kind of fiber. Mossy.—Organic soil material that contains more than 15 percent fibers that can be identified as moss and contains more moss fibers than any other kind of fiber. Herbaceous.—Organic soil material that contains more than 15 percent fibers that can be identified as herbaceous plants other than moss and grass or grasslike plants and more of these fibers than any other kind of fiber. In rare instances, some soil materials can be described by using a texture modifier, even though they do not fit the requirements of texture. An example is “gypsiferous material.” Limnic materials have modifiers to texture to connote the origin of the material. The three kinds of limnic materials are coprogenous earth, diatomaceous earth, and marl. These materials were deposited in water by precipitation or through the action of aquatic organisms or derived from plants and organisms. Refer to the Keys to Soil Taxonomy for the complete definitions and taxonomic criteria of limnic materials. The following three compositional texture modifiers are used with limnic materials to indicate presence and origin without respect to any set quantity of pellets, grains, or particles: Coprogenous.—Soil material that is a limnic layer containing many very small (0.1 to 0.001 mm) fecal pellets. Diatomaceous.—Soil material that is a limnic layer composed of diatoms. Marly.—Soil material that is a limnic layer that is light colored and reacts with HCl to evolve CO2. “Permanently frozen” is a texture modifier term applied to a soil layer in which the temperature is perennially at or below 0 degrees C, whether its consistence is very hard or loose. Entries.—Enter the applicable texture modifiers from the list in Part 618, Subpart B, Exhibits, Section 618.94. Multiple texture modifiers are used in some horizons based on the application rules for texture modifier presented above. They must be assigned sequence numbers in the Horizon Texture Modifier Table in the NASIS database for the proper calculated result. Texture Modifier and Class Definition.—“Texture modifier and class” is a concatenation of texture modifier and texture class or texture modifier and a term used in lieu of texture. This data element indicates the full texture term of the horizon. If texture modifiers are used, they are attached to the texture class by a hyphen, for example, GR-SL. If a layer is stratified, enter SR as a texture modifier and the end members of the textural range and connect them by hyphens, for example, SR-C-L and SR-GR-S-GR-C. Entries.—Enter the appropriate texture modifier and class for each horizon. These entries are calculated in the Horizon Texture Group Table in the NASIS database. Stratified Texture Flag Application.—A “stratified texture flag” is used to identify stratified textures in the Horizon Texture Group Table in the NASIS database. Entries.—A Boolean flag is set to “yes” by checking the box for the stratified texture flag. This indicates that the textures that comprise a particular record are stratified. The default entry is “no” and is displayed by keeping the box for stratified texture flag unchecked. Representative Indicator Flag Application.—A “representative indicator flag” is used to identify one representative texture (comprised of texture modifier and class) in the Horizon Texture Group Table in the NASIS database. Entries.—A Boolean flag is set to “yes” by checking the box for the representative indicator flag. This indicates that the texture that comprises a record in the particular horizon texture group is representative. It also indicates that the selected texture validates the soil properties populated for the layer. The selected texture record must be in agreement with the representative values for important soil properties such as clay content, sand content, rock fragment content, and organic matter content. The flag must be set even when only one texture record is populated for a particular horizon (such as in surface layers or bedrock layers). The default entry is “no” and is displayed by keeping the box for representative indicator flag unchecked. Only one texture record may be selected as representative for a given horizon or layer. 618.72 Water, One-Tenth Bar Definition.—“Water, one-tenth bar” is the amount of soil water retained at a tension of 1/10 bar (10 kPa), expressed as a percentage of the whole soil on a volumetric basis. Significance.—Water retained at one-tenth bar is significant in the determination of soil water-retention difference, which is used as the initial estimation of available water capacity for some soils. Measurement.—Measurement in the laboratory is done on natural clods using a pressure desorption method. Measurement for nonswelling soils, loamy sand or coarser soils, and some sandy loams is also done using a a pressure desorption method but sieved (< 2 mm) air-dry samples are used. Gravimetric water contents are reported in laboratory measurements as a percentage of the fine-earth (<2 mm) fraction. Conversion to a volumetric basis is made using bulk density and rock fragment content. Entries.—Enter the low, high, and representative values for the horizon. The range of valid entries is from 0 to 100 percent, and tenths (one decimal place) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.105. 618.73 Water, One-Third Bar Definition.—Water, one-third bar” is the amount of soil water retained at a tension of 1/3 bar (33 kPa), expressed as a percentage of the whole soil on a volumetric basis. Significance.—Water retained at one-third bar is significant in the determination of soil water-retention difference, which is used as the initial estimation of available water capacity for some soils. Measurement.—Measurement in the laboratory is done on natural clods using a pressure desorption method. Measurement for nonswelling soils, loamy sand or coarser soils, and some sandy loams is also done using a pressure desorption method but sieved (< 2 mm) air-dry samples are used. Gravimetric water contents are reported in laboratory measurements as a percentage of the fine-earth (< 2 mm) fraction. Conversion to a volumetric basis is made using bulk density and rock fragment content. Entries.—Enter the low, high, and representative values for the horizon. The range of valid entries is from 0 to 100 percent, and tenths (one decimal place) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.105. 618.74 Water, 15 Bar Definition.—“Water, 15 bar” is the amount of soil water retained at a tension of 15 bars (1500 kPa), expressed as a percentage of the whole soil on a volumetric basis. Significance.—Water retained at 15 bar is significant in the determination of soil water-retention difference, which is used as the initial estimation of available water capacity for some soils. Water retained at 15 bar is an estimation of the wilting point. Measurement.—Measurement in the laboratory is done on sieved (< 2 mm) air-dry samples using a pressure desorption method. Gravimetric water contents are reported in laboratory measurements as a percentage of the fine-earth (< 2 mm) fraction. Conversion to a volumetric basis is made using bulk density and rock fragment content. Entries.—Enter the low, high, and representative values for the horizon. The range of valid entries is from 0 to 100 percent, and tenths (one decimal place) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.105. 618.75 Water, Satiated Definition.—“Water, satiated” is the estimated volumetric soil water content at or near zero bar tension, expressed as a percentage of the whole soil. Significance.—Water, satiated, represents the total possible water content of the soil, including the amount in excess of field capacity, and is used to estimate the amount of water available for leaching and translocation. Satiated water content approximates the water content at saturated conditions. It is used in such resource assessment tools as Soil Hydrology, Water Budgets, Leaching, and Nutrient/Pesticide Loading models. Entries.—Enter the high, low, and representative values for the horizon. The range of valid entries is from 0 to 100 percent, and only whole numbers (integers) are allowed. A NASIS calculation is available and can be viewed in Part 618, Subpart B, Exhibits, Section 618.105. 618.76 Water Temperature Definition.—“Water temperature”, is the mean annual water temperature (MAWT) at or near the water/soil contact in a subaqueous soil setting. Significance.—Temperature is important to many biological and physical processes that occur in marine and freshwater aquatic environments. The properties of the water column above subaqueous soils are important to interpretations such as aquaculture, shellfish restoration, and seagrass survival. Estimates.—Water temperature can be estimated from measured water temperatures of the vicinity. Seasonal air temperature and water current fluctutations affect water temperatures. Estimates of water temperature should take these factors into account when water temperatures are extrapolated from one soil map unit component to another. Temperatures can be summations of the daily values collected and populated for point data in the NASIS data element “Water Temp – Lower” (sas_water_temp_lower) . Such temperatures are measured in the lower 10 cm of the water column immediately above the surface of subaqueous soils. Measurement.—Water temperature can be measured by many types of thermometers, including mercury, bimetallic, thermisters, and thermocouples. Many types of instruments can be configured for remote, unattended, and submerged operation. Entries.—Enter the high, low, and representative values for the range of mean annual water temperature for the component as the average of the mean monthly water temperatures in the Component table. The range of valid entries is from -10 to 50 degrees Celsius, and tenths (one decimal place) are allowed. 618.77 Wind Erodibility Group and Index Definition.—A wind erodibility group (WEG) is a grouping of soils that have similar properties affecting their resistance to soil blowing in cultivated areas. The groups indicate the susceptibility to blowing. The wind erodibility index (I), used in the wind erosion equation, is assigned using the wind erodibility groups. Significance.—There is a close correlation between soil blowing and the size and durability of surface clodiness, fragments, organic matter, and the calcareous reaction. The soil properties that are most important with respect to soil blowing are listed below. Soil moisture and the presence of frozen soil also influence soil blowing. Soil texture class Organic matter content Carbonates in the fine-earth fraction as determined by effervescence class Rock and pararock fragment content Mineralogy Estimates.— Soils are placed into wind erodibility groups on the basis of the properties of the soil surface layer. Subpart B, Exhibits, Section 618.95 lists the wind erodibility index assigned to the wind erodibility groups. The wind erodibility index values are assigned because the dry soil aggregates are very use-dependent on crop management factors. Entries.—Enter the wind erodibility group and wind erodibility index values for surface layers only. The range of valid entries for wind erodibility group data is 1, 2, 3, 4, 4L, 5, 6, 7, and 8. The lowest valid entry for wind erodibility index data is 0, and the highest is 310. The index values should correspond exactly to their wind erodibility group.

Atrial fibrillation (AF) is known as the most common arrhythmia in adults, and it is estimated that millions of people suffer from AF. Previous studies reported that AF may increase risk for other diseases such as stroke, cardiovascular mortality, and obstructive sleep apnea. Cardiac rhythms of patients with AF possess two particular characteristics: (i) R-R interval irregularity; and (ii) absence of P waves. Although the absence of P waves is a primary indicator of AF, P waves are low amplitude which makes them difficult to analyze due to high intensity noise in electrocardiograms (ECGs) derived from ambulatory electrocardiography. Compared to P waves, R waves have much higher amplitudes and therefore higher tolerance to ECG noise, making detection mush easier. Consequently, using R-R intervals to detect AF is considered to be the most robust approach. Accordingly, several AF detection methods based on R-R interval irregularity have been proposed. However, the R-R interval window length required for those methods varies from 32 beats to 128 beats. To exhibit high accuracy and finer temporal resolution for real-time monitoring and shorter AF episodes detection, computational complexity together with the power consumption is dramatically increased.

Immunizations. Vaccination is one of the most effective preventative health measures in modern times. It has positively affected the incidence and prevalence of a myriad of diseases, as well as improved both their morbidity and mortality. It is essential that health care providers understand the value of this intervention, promote comprehensive immunization coverage within their community, and recognize the underlying principles that guide decisions regarding vaccination. Factors important to consider before immunization administration include the type of vaccination, its recommended timing and intervals, and the presence of contraindications, precautions, or other special considerations.

Sure, you’re trapped in your house, but there’s an entire galaxy of cool stuff right over your head in the night sky. Space is infinitely expanding (and contracting, but we’ll focus on the expanding), which might relieve the claustrophobia and boredom that come along with quarantine. Here are eight pictures showing the mystery and wonder of the universe. Star Cluster Just how it sounds, a star cluster is a bound group of stars. Young clusters contain fewer stars with more distance between them. As a cluster gets older, the stars move closer together, united by gravity. You can see some star clusters in the night sky with the naked eye. One is the constellation Seven Sisters, another’s inside the constellation Cancer. The closest star cluster, Hyades, is just 153 light-years from earth and you can see it inside the constellation Taurus. Nebula The word nebula is Latin for cloud, which is also a perfect description. Images of nebulae from space are beautiful and other-wordly. They occur when dust, hydrogen, or ionized gases come together to form a cloud in space. There’s one visible to the naked-eye, the Orion Nebula. It’s positioned inside the Orion constellation, just below the belt of three stars. Mars Doesn’t this look like it could be a picture of deserts on Earth? This one is from the Mars Reconnaissance Orbiter, taken back in 2015. It shows the dunes in a desert on the surface of the planet. (x) Andromeda Galaxy Our nearest neighboring galaxy, just 2.5 million light years away. Early astronomers first observed it in the 10th century, calling it a nebulous smear. Continued study through the centuries led to the first observation of a supernova in the 19th century. Not until 1925 was it identified as a separate galaxy by Edwin Hubbard, using his Hooker telescope. Aurora Borealis When solar wind hits the upper atmosphere of earth, it changes the paths of particles, creating beautiful light shows visible both from the surface of the planet and outer space. Sirenum Fossae, Mars There are several long troughs on the surface of Mars, called fossae. Gullies appear throughout the Sirenum Fossae, formed by frozen carbon dioxide running like water down the faces of dunes and craters. Total Solar Eclipse In 2017, the northern hemisphere of Earth observed a total solar eclipse, which happens when the orbit of the moon brings the Earth into perfect alignment with the Sun and Moon. There’s between 2 and 5 eclipses every year, with an average of 240 per century. The next time there will be 5 eclipses in a single year will be in 2206. Triangulum Galaxy The most distant permanent celestial body that’s visible to the naked eye, it’s smaller than both the Milky Way and Andromeda. It’s moving toward the Andromeda galaxy, leading astronomers to believe it’s a satellite of the larger galaxy.

Quadrilateral ABCD has coordinates A( 3, -5) B (5, -2) C (10, -4) , D (8, -7) Quadrilateral ABCD is a rectangle, because opposite sides are congruent and adjacent sides are perpendicular Square, because all four sides are congruent and adjacent sides are perpendicular parallelogram, because opposite sides are congruent and adjacent sides are not perpendicular rhombus, because all four sides are congruent and adjacent sides are not perpendicular OpenStudy is now Brainly! Quadrilateral ABCD has coordinates A( 3, -5) B (5, -2) C (10, -4) , D (8, -7) Quadrilateral ABCD is a rectangle, because opposite sides are congruent and adjacent sides are perpendicular Square, because all four sides are congruent and adjacent sides are perpendicular parallelogram, because opposite sides are congruent and adjacent sides are not perpendicular rhombus, because all four sides are congruent and adjacent sides are not perpendicular Mathematics Stacey Warren - Expert brainly.com Hey! We 've verified this expert answer for you, click below to unlock the details :)

Pediatric surgery between 1860 and 1900. The twentieth century begins the modern era of pediatric surgery, a time of innovation and vastly improved operative survival. Until then, many correctable congenital anomalies had been described in case reports only as medical curiosities, seemingly having no chance for surgical cure, but the challenge was available to those wishing to accept. The first texts devoted to surgical diseases of infancy and childhood were descriptive, emphasizing visible and palpable lesions, tuberculosis, musculoskeletal deformity, and major trauma, but also recognized that "surgical diseases of children by their number, their variety, and their character must, we think, be studied in a special manner and demand a special chapter in surgical books." The construction of Children's Hospitals during the nineteenth century identified the need to provide special facilities for children. Surgical antisepsis was practiced and general anesthesia used. Physicians such as Hirschsprung recognized a wide spectrum of treatable surgical conditions that would soon become the domain of pediatric surgeons. Indeed, at that moment in time at the very beginning of the twentieth century, Roentgen discovered the x-ray, Ladd received his MD degree at Harvard, the first pediatric appendectomies for appendicitis were successfully performed, and Fredet and Ramstedt recognized that pyloroplasty was not the preferred treatment for hypertrophic pyloric stenosis. Modern pediatric surgery had passed from embryo to fetus to newborn.

Abstract The dynamics of a wild red-legged partridge population were examined over a 14-year period in Spain to identify patterns in age and sex ratios in relation to weather parameters, and to assess the importance of these parameters in population dynamics and management. The results gave age ratios of 1.07 (but 2.13 in July counts), juvenile sex ratios of 1.01 and adult sex ratios of 1.47. Overall, 12% more females were hatched and female juvenile mortality was 7.3% higher than in males. Sex differential mortality explains the 19.2% deficit in adult females, which are more heavily predated than males during the breeding period. Accordingly, age ratios are dependent on sex ratios and both are density dependent. Over time, ratios and density changes appear to be influenced by weather and management. When the habitat is well conserved, partridge population dynamics can be explained by a causal chain: weather operates on net primary production, thereby affecting partridge reproduction and predation and, as a result, age and sex ratios in the October population. A reduction in the impact of predation (i.e. the effects of ground predators on eggs, chicks and breeding females) is the key factor to improve the conservation of partridge populations and associated biological processes. Citation: Nadal J, Ponz C, Margalida A (2016) Age and Sex Ratios in a High-Density Wild Red-Legged Partridge Population. PLoS ONE 11(8): e0159765. https://doi.org/10.1371/journal.pone.0159765 Editor: Ulrike Gertrud Munderloh, University of Minnesota, UNITED STATES Received: March 10, 2016; Accepted: July 6, 2016; Published: August 10, 2016 Copyright: © 2016 Nadal et al. This is an open access article distributed under the terms of the Creative Commons Attribution License, which permits unrestricted use, distribution, and reproduction in any medium, provided the original author and source are credited. Data Availability: The minimal data necessary to replicate the study findings has been uploaded as Supporting Information. Funding: AM was supported by a Ramón y Cajal research contract from the Ministry of Economy and Competitiveness (RYC-2012-11867). Competing interests: The authors have declared that no competing interests exist. Introduction Age and sex ratios change as annual partridge production and predation rates vary [1–3] and hence large age ratios indicate good annual partridge recruitment. Juvenile sex ratios are generally 1:1; however, adult sex ratios in birds are generally biased towards males but in mammals are frequently skewed towards females [4–6]. Different genetic, behavioural and ecological theories have been formulated to explain biased sex ratios during life stages [5–7], and characteristics ranging from genetic sex determination to parental condition, behaviour, environment and weather have all been postulated as the main factors capable of modulating juvenile sex ratios [8–12]. Age and sex ratios in partridge populations are still very poorly known [13–15]. The frequency and causes of variations in these ratios and the responses in population dynamics are all poorly understood and so the overall extent to which subtle changes in age and sex ratios can modify partridge densities is also unknown [13,16,17]. Understanding variations in age and sex ratios is extremely challenging because ratios vary seasonally, annually, spatially and between age and sex classes, and require clear definitions of what constitutes a sample: for instance, samples from smaller or larger spatial areas, or over shorter or longer time intervals, are interpretable in different ways [18–20]. The development of a simple and robust method for estimating age and sex ratios in the field, as well as for interpreting the results obtained, could improve wildlife management and conservation. Habitat quality (resource availability), density (competition), sexual size dimorphism (sex specific mortality), mating system (sexual selection, parental condition and capacity), synchrony (distribution of temporal opportunity) and predation (population regulation) affect offspring condition and production (chick numbers, quality and sex ratio) [4,6,21]. A better understanding of age and sex ratios in wild game bird populations could greatly improve our understanding of population processes and contribute to their management and conservation [22–24] The increase in anthropogenic activities is having an impact on natural ecosystems and degrading biological processes [25,26]. Consequently, wildlife management has had to urgently establish conservation priorities as the degradation of ecosystems derives towards a loss of natural resources. In recent years, wild partridge populations have been replaced by farmed birds and game workers (managers, keepers, agents and businesses) have shifted their activities towards artificial production. In red-legged partridges (Alectoris rufa), rapid habitat degradation has led to increased exposure of nests and coveys to predators and detrimental anthropogenic activities (i.e. agricultural machinery). This provokes greater mortality in young and female birds and thus consistently adult-skewed and male-skewed age and sex ratios in these populations [27]. Hence, in areas of better quality habitat, we can expect juvenile survival and age ratios to be higher, and also that male and female survival rates will be similar and sex ratios balanced [3,28,29]. In this sense, net primary production (NPP) depends on meteorological conditions but is also associated with the availability of food resources and predation risks. Both factors alter age and sex ratios and as consequence, their dynamics reflect a population responses to changes in annual habitat quality [3,30,31]. Annual, local and global age and sex ratios in declining red-legged partridge populations have been reported in northern Spain [14,27,32]. By contrast, age and sex ratios in stable or expanding populations are still very poorly studied and long time datasets are costly and difficult to obtain. Monitoring game bird populations requires robust methods for annual assessments of fluctuations in populations and in the resulting population dynamics [31,33,34]. Here, we take advantage of data from a long-term study (14 years) of a wild red-legged partridge population in La Mancha (Spain) to (1) describe current age and sex ratios in a high density population; (2) identify patterns of age and sex ratios and their causes, correlates and variations in order to respond the following questions: (a) Are age and sex ratio mutually dependent? (b) Are age and sex ratio density dependent? (c) Are age and sex ratio of year n+1 dependent on age and sex ratios of year n?; (3) assess the importance of known abundance, age and sex ratios in population monitoring, management and conservation to evaluate population health, stage and viability. (4) use and propose various statistical techniques for correctly interpreting annual age and sex ratios in high density populations; and (5) examine the effect of weather on age and sex ratios. Materials and Methods Ethical statement The study was conducted in full compliance with Spanish laws and regulation, including the licence of “Las Ensanchas” to sampling shoot partridges. The protocol was approved by the Committee on the Ethics of Animal Experiments of the University of Lleida (Ref.1998-2012/05). Study area The study was carried out in Las Ensanchas, a small-game hunting estate in the basin of the river Jabalón in Ciudad Real province (central Spain; 38°39’ N, 3°13’ W; 790–840 m a.s.l.; annual average temperature is 13° degree Celsius (range 4–24, January and August, respectively) and annual average precipitation is 371.8 mm (range 5–725, July-December, respectively). It covers a surface area of 1,415 ha of predominantly species-rich Mediterranean landscapes (wood pastures or dehesa) that are characterised by a habitat mosaic consisting of cereal crops, fallow, natural pastures (Agrostis castellana, Poa bulbosa, Stipa tenacissima, Brachypodium distachyon, Trifolium glomeratum, Bellis annua, Bellis perennis) and scrubland (Cistus ladanifer, Quercus coccifera, Thymus vulgaris, Rosmarinus officinalis, Lavandula latifolia, Salvia officinalis, Santolina chamaecyparissus, Helianthemum syriacum) with scattered holm oaks (Quercus ilex). In all, 75% of the estate is covered by herbaceous vegetation and the other 25% by shrubland. The main quarry species are rabbit (Oryctolagus cuniculus), red-legged partridge (Alectoris rufa), woodpigeon (Columba palumbus) and hare (Lepus granatensis). Controls of predators annually remove on average 40 red foxes (Vulpes vulpes), 70 feral cats (Felix domesticus) and 300 magpies (Pica pica). This property also holds more than 150 important and protected bird species including Spanish imperial (Aquila adalberti) and golden (Aquila chrysaetos) eagles, cinereous vulture (Aegypius monachus), pin-tailed (Pterocles alchata) and black-bellied (Pterocles orientalis) sandgrouse, among others. Data collection We analysed red-legged partridge hunting bags from the period 1998–2011. In the field, we examined all bagged birds, certified and classified them as wild partridges. The annual average hunting yield was 0.77 partridges/ha, 33% of the autumn population (2.31 partridges/ha). From spring (1.14 partridges/ha) to autumn, the population increased 51% (Supporting Information, S4 Appendix). Hunting methods (drives), as well as the ability of hunters and beaters, remained constant over the years. The team at Las Ensanchas worked for profit as all shot partridges were sold as a high-quality food product. A total of three to six people processed the warm, recently bagged birds. Age was determined by an examination of primary feathers and sex by spur characteristics [35–39]. Clear reference patterns from a previously identified bird for each age and sex class (young or adult; female or male) were always available. A wing was taken from all birds for precise measurement in the laboratory to confirm the age and sex determination carried out in the field and thus to verify the class they had been assigned to [37,38]. Statistical analysis We used age ratio (young/adults) to describe the population age structure, sex ratio (male/female) to describe the population sex structure, adult sex ratio (adult male/adult female) to describe the adult sex structure, and juvenile sex ratio (young male/young female) to describe the juvenile sex structure. We examined scatterplots to gain a thorough understanding of the constraining assumptions imposed by our data set [40]. The lack of data independence stemmed from a homogeneous biological structure, as within the population there is a strong family relationship between individuals. Many partridges in the sample came from the same family or the same lineage. These specimens could be more similar to each other than to individuals from other locations or years that are separated by larger distances or longer time periods. We associated autumn ratios with the summer age ratios derived from counts conducted in July [15]. The two indexes are very different because of different collection methods. We tested for the biological significance of the results and initially applied different procedures and statistical tests to confirm results and to make robust interpretations. The most robust tests were maintained to simplify the results obtained. To test ratio independence, we used simple regression analyses. For the response, we used ratios resulting from the 1) integer fraction, 2) proportion = dividend / dividend + divisor, and 3) logarithm fraction. We employed frequency = count individuals of one type; absolute frequency = total sample number individuals; relative frequency = proportion; and class (groups of birds according to age and sex) [41]. We used the Chi-square test to determine ratios that differed from a 1:1 ratio as a reference, and contingency tables to test the null hypotheses that sexes and/or ages were distributed randomly over the years. We employed simple and multiple regression analyses for the response variables (ratios, proportions, class) explained by other ratios, proportions, class, year, and other independent variables such as density or weather. In addition, to fit age and sex ratio responses, we used (a) logistic regression, with model terms: absolute frequency, sample size, year, density and complementary age or sex; (b) generalized regression models, allowing small data sets and multicollinearity, with different distributions to find the best fit; (c) generalized linear models (GLMs), with different numbers of parameters; and (d) mixed models, which take into account complex covariance structures[42]. All tests were considered significant when P < 0.05 or when similar comparisons were significant and supported tests with P < 0.1. We applied principles of parsimony for rationalising the model and only variables that contributed significantly were retained. The model did not contain any redundant parameters or factor levels and the number of parameters and data were balanced [41]. The AICc (corrected Akaike information criterion), delta AICc and the Akaike weights were used to assess different models, as well the significance level of models and effects [43]. In the GLMs, we applied the deviance as a measure of the goodness of fit. To build the models, we searched for all factors, covariates and interaction terms that might influence ecological explanations. We then carried out a series of step-wise deletions of any non-significant explanatory variables, factors and interaction terms, and continued to remove terms until the best model was obtained [40,43]. To evaluate the effects of weather on age and sex ratios we used data from Las Terceras, a meteorological station run by the AEMET (Spanish Meteorological Agency, http://www.aemet.es) located near the edge of the estate (300 m). We selected the best rain and temperature data according to the relationship between vegetation growth and partridge parameters, and selected the months and variables that showed the greatest correlation with the partridge population [43]. Discussion Sex ratios, chick survival (the key factor in population size), excess males and high chick mortality are all related to population density [2,44,45]. Our findings suggest that there was no significant relationship in ratios between consecutive years [17,22,46]. Density was explained by age (31% of variance) and sex ratios (32% variance) [45], although other factors did have some effect [33,47,48]. Age ratios were dependent on sex ratios (32% of variance) [19,49] and reflected the results of annual productivity minus mortality in the data-collection period. High annual brood production and age ratio variation may mask density-dependent responses [34,50,51]. Abundance partially depends on age and sex ratios [52,53] and in high density population, partridge interactions are more likely and influential. Density depends partially on population structure and age ratios partially depend on sex ratios. The population structure is partially intra-dependent. Inter- and intra-dependence of population structure are partial because these associations can be mediated by other factors (e.g., predators, availability of food and cover, diseases and weather) [54–58]. Predation was responsible for a decrease in the age ratio by half between July and October [1,15,51,55]. Our results showed that 12% more females hatched and had 7.3%-higher juvenile mortality than males, a sexually differentiated mortality that continues into adult age classes. There were 19.2% fewer adult females than adult males, and so the adult sex ratio was male-biased [46,55,59]. The characteristics of the population structure begin to form in the growth stages of partridges (from hatching to juvenile), although the process continues in the maturation phase to adult (from juvenile to adult), with predation being the principal driver of population structure [55,58,60]. From a methodological point of view, all the statistical tools applied were useful for understanding age and sex ratios and the temporal changes they undergo, as well as their relationships with other factors (effects) that contribute to explaining variation. The different statistical tools generate complementary approximations for explaining age and sex ratios. The most complete models explain many different effects, while the simplest show the degree of association [61,62]. For example, our results indicated that multiple regression demonstrates that meteorological parameters explain age and sex ratios, while simple regression illustrates how each meteorological parameter explains these ratios. Age and sex ratio are associated with spring rainfall and temperature. Accordingly, statistical techniques condition our biological interpretation of results and so we must consider how they influence our understanding of wildlife populations to make wide-ranging interpretations. For example, logistic regression for explaining age ratio supports the idea that year and sex had a significant effect, as did the sample size for sex, but density was not included in the effects. Generalized regression models explained the relative frequency with density, year and class effects, while generalized linear models using binomial distribution explained class, sex and age, with effects for absolute frequency, sample size, year and density (Supporting Information S2 and S3 Appendices). When we manage high-density partridge populations average production and predation maintain close to a 1:1 age ratio in the bag (October). We anticipate that, under future scenarios of adverse meteorological conditions, the NPP will be poorer and will entail a fall in density and age ratios. By contrast, when weather conditions are more favourable, the NPP, density and age ratios will all rise. Partridge populations at high densities maintained a balanced age ratio of around 1:1. Age ratios differed from the period mean (1.07) in 12 of the 14 studied years, thereby indicating that variability in years with lower or higher productivity rates is very frequent. The annual extraction rate should be adjusted to the age ratio of that year [15,63,64]. Sex ratios showed little change over the years, the only variation being change in juvenile sex ratios in some years [33,34]. Accordingly, it is difficult to identify the direct cause behind this population response in the juvenile sex ratio, and it would appear to be the cumulative result of several effects: sex differential in births, mortality, resource availability, predation risk, local storms, local heat waves and other possible causes [9,10,59,65]. Sex ratios only differed from the period mean in two years. Although the juvenile sex ratio showed a little annual change, the adult sex ratio was constant over the entire period and corresponds to the maximum partridge density in a managed habitat, where changes in the availability of resources depend on weather conditions [3,30,66]. In spite of efficient predator control by gamekeepers, adult sex ratios indicate a differential loss of females. Adult females are more exposed to predation than males as a consequence of their parental effort during incubation and chick rearing [18,21,67]. Adult sex ratio distortion in birds is significantly more severe in populations of globally threatened species than in non-threatened species [13,20,23], and greater female mortality is the main driver of male-skewed adult sex ratios in birds [68,69]. Thus, predator control and habitat improvement are needed to balance sex ratios. We can assume that age ratio homeostasis over decades is the response of a population in a stable habitat where density oscillations occur due to variations in NPP. In high-density partridge populations, longevity increases and production is checked due to higher levels of competition and predation [25,28,53,70]. The same values of age and sex ratios have different population effects in low, middle and high density populations. Consequently, monitoring age and sex ratio controlling for density is necessary to manage and conserve wildlife populations. An useful method for understanding population responses to ecosystem management is required and, accordingly, our study shows how age and sex ratio dynamics are key tools for interpreting, managing and conserving wildlife populations. Supporting Information S1 Appendix. Determination coefficient (R2) and slopes (b±SD) in regression between abundance (Y) and the age and sex ratios (X). Determination coefficient (R2) and slopes (b±SD) in regression between age ratio (Y) and sex ratio (X). Determination coefficient (R2) and slopes (b±SD) in regression between age and sex ratio and the previous year age and sex ratio. Effects, corrected Akaike information criterion (AICc), with logistic regression models for age or sex explained by absolute frequency, sample size, year, density, sex or age. https://doi.org/10.1371/journal.pone.0159765.s001 (DOCX) S2 Appendix. Generalized regression models with relative frequency explained by class density, year and class. Generalized regression models with absolute frequency explained by number of trials, class density, year (ordinal), age and sex. Generalized regression models with absolute frequency explained by number of trials, age, sex, year and density. Generalized regression models with absolute frequency explained by number of trials, age, sex, class density and year (ordinal). https://doi.org/10.1371/journal.pone.0159765.s002 (DOCX) S3 Appendix. Generalized linear models (GLMs). Relative frequency, age, sex, year and density. Relative frequency, class, year and density. Absolute frequency, trial, age, sex, year and density. Absolute frequency, trial, class, year and density. https://doi.org/10.1371/journal.pone.0159765.s003 (DOCX) S4 Appendix. Population numbers. https://doi.org/10.1371/journal.pone.0159765.s004 (DOCX) Acknowledgments We are grateful for contributions made by the Melgarejo family, Patricia, Luis and Ivan Maldonado and Tom Gullick. Thanks also go to the “Las Ensanchas” staff, especially the game keepers, the Barranquero family and collaborators, the members of the Tom Gullick hunting team in “campo de Montiel” and around the world. Alberto Meriggi, Dick Potts and Christian Gortazar provided valuable comments and significantly improved the manuscript. Author Contributions Conceived and designed the experiments: JN. Performed the experiments: JN CP. Analyzed the data: JN CP. Contributed reagents/materials/analysis tools: JN CP. Wrote the paper: JN AM.

Q: Recursion: look like same function but print out different flow of execution I tried to print out flow of execution of factorial recursion function. def factorial(n): space=' '*(4*n) print(space,'factorial ',n) if n==0: print(space,'returning 1') return 1 else: recurse=factorial(n-1) result=n*recurse print(space,'returning ', result) return n*recurse With factorial(2) result will be correct: factorial 2 factorial 1 factorial 0 returning 1 returning 1 returning 2 2 But when I don't use variable recurse=factorial(n-1): def factorial(n): space=' '*(4*n) print(space,'factorial ',n) if n==0: print(space,'returning 1') return 1 else: result=n*factorial(n-1) print(space,'returning ', result) return n*factorial(n-1) then with factorial(2) it printed out different flow: factorial 2 factorial 1 factorial 0 returning 1 returning 1 factorial 0 returning 1 returning 2 factorial 1 factorial 0 returning 1 returning 1 factorial 0 returning 1 2 My question is: why was there different flow of executions [with and without variable recurse=factorial(n-1)]? A: This is due to below extra function call in your code def factorial(n): space=' '*(4*n) print(space,'factorial ',n) if n==0: print(space,'returning 1') return 1 else: result=n*factorial(n-1) print(space,'returning ', result) return n*factorial(n-1) # factorial(n-1) is called second time so that this function is executed once again You could do as below for similar execution def factorial(n): space=' '*(4*n) print(space,'factorial ',n) if n==0: print(space,'returning 1') return 1 else: result=n*factorial(n-1) print(space,'returning ', result) return result

Cameras of all types are becoming more and more affordable to consumers as technologies for manufacturing cameras have become more efficient and cost effective. Consequently, many consumers are purchasing and using cameras to take photographic pictures, video recordings, and even digitized pictures. Many amateur photographers or videographers make mistakes when taking pictures or making video recordings which can only be avoided by taking photography or videography courses. One such area of common mistake is in the area of picture composition. Because people have an affinity for symmetry, amateur photographers or videographers tend to align a subject in the center of a picture or video recording at all times, creating bland or sometimes disastrous picture composition. These amateurish mistakes are further reinforced by cameras currently in use. As shown in FIG. 1, viewfinders of most cameras only have a circle or a cross-hair in the center of a viewfinder with maybe brackets on the outer edges to delimit the size of pictures or screens on which the picture would be transferred. Almost instinctively, users align the subject of the picture to the circle or cross-hair in the center of the viewfinder. Even those users who are taking or have taken picture composition courses sometimes do not understand the concept of horizons, diagonals, and strong points, for example, or forget the lessons learned from lack of frequently practice. What is needed is a user-friendly guidance system that would aid a camera user in taking pictures with quality picture compositions.

Everything about Bacteriophage Menu Lambda Phages The lambda phage, also called Enterobacteria phageλ andcolphageλ, is a type of temperate bacteriophage or bacterial virus that infects the Escherichia coli (E. coli) species of bacteria. The virus may be housed in the genome of its host via lysogeny. History of Lambda Phage In 1950, Esther Lederberg, an American microbiologist, was performing experiments on E. coli mixtures. She happened to observe streaks of mixtures of two types of E. coli strains that seemed to have been nibbled on and had viral plaque. One E. coli strain had been treated with ultraviolet light, so the damage was more apparent. It was later determined that this was caused by bacterial viruses, which replicated and spread resulting in cell destruction. The discovery led to the employment of Lambda phage as a model organism in microbial genetics as well as in molecular genetics. Structure A lambda phage has a head measuring around 50-60 nanometers in diameter and a flexible tail that is around 150 nanometers long and may contain tail fibers. The head consists of various proteins and over a thousand protein molecules including X1, X2, B, B*, E, D, and W. The head functions as a capsid that contains its genome, which contains 48,490 base pairs of double-stranded linear DNA. This number also includes 12-base single-stranded parts at its 5’ ends. The single-stranded parts are known as sticky sites and are also called the cos site, which encircles the DNA in the host cytoplasm. Hence, when in circular form, the phage genome is comprised of 48,502 pairs in length. The weight of the genome is estimated to be 32×106 Da, which is around half of the weight of the phage. The tail has a 135 nanometer-tube that is hollow and contains a conical cap which is around 15 nanometers. The tail’s inner diameter is around 3 nanometers, while on the outside, it is around 9-18 nanometers depending on the knob-like structures that give the tail a rough appearance. Life cycle When E. coli is infected with a lambda phage, two cycles may happen: lytic or lysogenic. The lytic cycle happens when progeny phage particles are produced. The lytic cycle is the more common life cycle that comes after most infections. The first step of this cycle is the attachment of the phage to the host cell, injecting its DNA into the cell. Nucleic acid from the phage is replicated, and the phage’s genes are expressed, allowing the production of phage proteins. The phage proteins are assembled into phage particles, which are released when the host cell undergoes lysis (it breaks down). The lysis is mediated by lysis genes S, R, Rz, and Rz1 which, upon expression, yield proteins that work together to break down the host bacterium’s cell wall. This mode of lambda replication typically yields many phage particles. The lysogenic cycle, in contrast, does not produce a huge number of progeny phage or break down the host cell. Instead, the λ DNA recombines with its host’s genome to produce a prophage. This typically is the favoured pathway when unfavourable environmental conditions prevent intense replication of the bacterial cells. Like the lytic cycle, the first step of the lysogenic cycle is also the attachment of the phage and the injection of its DNA into the host cell. The phage DNA then integrates with the host chromosome, producing an integrated DNA combination called the prophage DNA. Host cells that carry this DNA are said to be in the lysogenic state. The prophage DNA is replicated along each time the host bacterial cell replicates itself, producing more cells, each with a copy of the prophage DNA. When these cells are exposed to certain chemicals or to ultraviolet light, phage induction happens; the prophage DNA is then cut out of the host genome and proceeds to the lytic cycle. Applications The lambda phage has different applications, most of which are related to DNA cloning. This is because lambda phage can be used as a vector for generating recombinant DNA, which are combined DNA sequences that result from using laboratory techniques like molecular cloning to assemble genetic material from several sources. The site-specific recombinase of lambda phage can be used for shuffling cloned DNAs via the gateway cloning system, a molecular biology technique that ensures the effective transfer of DNA fragments between plasmids. The lambda phage’s ability to mediate genetic recombincation is due to its red operon, which is a functioning unit of genomic DNA that has a cluster of genes controlled by a promoter or a single regulatory signal. This red operon can be expressed to yield the proteins red alpha (or exo), beta, and gamma, which can be used in recombination-mediated genetic engineering, a method commonly employed in bacterial genetics, generation of target vectors, and DNA modification.

The disclosure relates generally to hot gas path (HGP) components, and more particularly, to a cooling structure for an HGP component and methods of fabricating the cooling structure. In a turbine system, a “hot gas path” (HGP) component can generally include a nozzle, rotor blade, shroud, or other hardware used in a combustion section of a gas turbine. Stationary blades are a type of HGP component used in turbine applications to direct hot gas flows to moving rotor blades, also known as buckets, to generate power. In steam and gas turbine applications, the stationary blades are referred to as nozzles, and are mounted to an exterior structure such as a casing and/or an internal seal structure by endwalls. Each endwall is joined to a corresponding end of an airfoil of the stationary blade. Stationary blades can also include passages or other features for circulating cooling fluids which absorb heat from operative components of the turbomachine. The flow of operative fluids, e.g., hot gas, against the buckets can cause a shaft to rotate and thereby drive a generator operatively connected to the turbine section of the system. In order to operate in extreme temperature settings, HGP components such as moving and stationary blades need to be cooled. In an example system, cooling fluids can be pulled from the wheel space or a dedicated source and routed into or through a cooling circuit of the blade to provide heat transfer between the composition of the blade and the cooling fluid. In contrast, in many gas turbine applications, later stage nozzles may be fed cooling fluid, e.g., air, extracted from a compressor of the gas turbine. In addition to the effectiveness of cooling, the structure of an HGP and cooling structures therein can affect other factors such as manufacturability, ease of inspection, and the durability of a turbomachine.

‘Then you lead me to the safe cave, and show / Me there my self, and in my own breast / Secret, deep miracles reveal themselves.’ Faust, Goethe INTRODUCTION: THE PHILOSOPHY OF KANT Modern philosophy began with René Descartes. He was the first person in the Western tradition to question the given ‘objective’ world. He was a complete lunatic, who believed that thought and identity were the same thing, that the physical realm and the mental realm were distinct components and that animals were machines, but we have him to thank for turning our attentions inwards and for taking the first faltering steps on a path of ‘radical scepticism’ which, in the end, leads to a miraculous — and very real — destination. For a century and a half after Descartes, philosophers swung between idealism — the belief that only the subjective mind is real — and realism — the belief that only the objective world is real. Most people today, most of the time, still swing between these two beliefs. They go by different names or, in the case of ‘the man on the street’, by no name at all, but idealism and realism roughly mark out the two poles of intellectual ignorance. They might still mark them out were it not for the appearance of Immanuel Kant who, in 1781, gave clumsy and confusing birth to the radical truth; transcendental idealism. Kant was an intellectual genius who had very little understanding of anything beyond the intellect, and so it is very hard to see how he reached the epochal conclusion that he did; that while the intellect can know facts, it can never know the truth; that reality is always, ultimately, beyond the reach of mind. Kant argued that time and space must be functions of the mind, which is to say, that the intellect has inherent limits, which can never be overcome. Reading his books, or trying to, and looking over some of the events of his life, as well as the famous revisions he made to the second edition of his most famous work, suggests that he came to these conclusions despite himself, that he was unwilling to accept that the entire manifest universe is a creation of the mind. But in any case, that is what he demonstrated. He pointed out that there is no way that the world out there can ‘migrate into my powers of representation’. That there is no way we could know anything of the world, as it exists in space and time, unless we already fundamentally knew space and time. There is, as Kant’s predecessors, Locke, Berkeley and Hume had said, just no conceivable way we could learn space and time (causality, matter, induction, etc, etc) from experience. Where Kant differed from these famous sceptics was in showing that we are born with a space-time faculty which can give us certain knowledge, but only of appearances, never of reality. ‘…if one eliminates from the empirical intuitions of bodies and their alterations (motion) everything empirical, that is, that which belongs to sensation, then space and time still remain, which are therefore pure intuitions that underlie a priori the empirical intuitions, and for that reason can never themselves be eliminated.’ (Prolegomena to Any Future Metaphysics). This means that you can ‘think away’ all the relative qualities of an object — its position, its colour, its shape, its state, etc. — because these are learnt, they come to us from acquired sensory experience; but you can never think away time and space — you can never perceive or conceive of [an object in] ‘no space’ or ‘no time’ — because you are born with this understanding ‘hardwired’ into perception and conception. Kant argued that our experience of space and time out there would be completely impossible unless we brought space and time to our experience and that there is no way to even imagine how it could be otherwise. Prior to Kant, John Locke had shown that the ‘secondary qualities’ of objects — their colours, smells, temperature, etc — were created by the [mental] organs of sense, and therefore could not exist in any meaningful way ‘out there’, in the world. Of course they do somehow exist out there — it would be utterly absurd to say that the sky is in no sense blue — but not in a way we can actually grasp, because all we experience is what our mind — through our senses — presents to us. Again any other way of understanding the matter is impossible to imagine or meaningfully talk of. Kant, following Berkeley, went on to demonstrate that ‘primary qualities’ are also mind-made; form, number, duration through time and so on do not really exist out there. Our minds create for us the seperateness of things, their location, their coming into existence and their passing away. Kant says that it is as absurd to assume that the appearance of the world resembles the world as it is to assume that ‘the sensation of red is similar to the property of cinnabar that excites this sensation in me.’ Again this does not mean that there is nothing in what we call a ‘table’ that makes it round, rather than square. Kant is not an idealist, or solipsist — he doesn’t believe that our minds literally make up the world from nothing (or God’s mind; as Berkeley did). Kant says there is a table out there, it is in some way square, but what that really means, what it really is, I can never directly experience. I can only know the table’s representation not the thing-in-itself. Kant therefore made an entirely different distinction between objective and subjective to that which preceded and followed him. Today we all more or less agree with Descartes that the objective world is ‘out there’ and the subjective world is ‘in here’. For Kant they are both ‘in here’, both elements of the same subject — or, looked at another way, both objects to a consciousness which precedes them. The problem for Kant is what that consciousness is. As he had no experience of it, his reasoning about it failed. He knew it was there, but that’s all he did — know. This is why he said, quite rightly, that nothing can be abstractly known of ‘I’, the psyche (i.e. that psychology is bullshit), but then assumed that it was thereby existentially unknowable — that I could not be experienced. This of course makes excellent sense to academics, who have almost zero conscious experience, but to anyone actually alive, rather than merely living, is complete nonsense. Kant’s hyper-abstraction also led to him rowing rapidly back from the marvellous discovery he had made, and in the second volume of his magnum opus, declaring, ludicrously, that the foundation for morality was conceptual. Everyone who has studied modern philosophy knows Kant’s conceptual formula for the good or the right; the fatuous bromide ‘act only according to that maxim whereby you can, at the same time, will that it should become a universal law,’ aka ‘the categorical imperative.’ By pulling the idea of moral goodness out of thin air Kant effectively paved the way for a veritable landslide of bullshit, culminating in the ‘hollowest verbiage’ of the ‘insolent nonsense-scribbler’ Hegel, who did nothing but pull stupendous sounding guff out of his arse. If reality is the rational mind, then why not? This ‘reality’ Hegel falsely, and absurdly, called ‘consciousness’ or ‘spirit’ (geist), a quasi-mystical — yet entirely rational and knowable — principle which basically argued with itself throughout history until it proved that freedom is rationality, unmediated experience of reality is a lie and the verbally incommunicable is by its very nature untrue; oh, and the German police state is paradise. That’s enough on Hegel’s nonsense though, which doesn’t deserve much serious consideration, and on Kant, whose most important arguments, those that can be understood without his ‘complicated clockwork,’ are better approached through the perceptive clarity of Arthur Schopenhauer. THE PHILOSOPHY OF SCHOPENHAUER Of all the offence Schopenhauer has given to numerous scholars, nothing has offended them more than the unfortunate fact that he does not resemble them. Untimely Meditations, Friedrich Nietzsche The paragraph above criticising Kant’s moral philosophy and Hegel’s chicanery came from the next — and some might say the last — Western philosopher to appear with anything meaningful to say in this tradition, and on this subject; the great Arthur Schopenhauer. Schopenhauer understood that although the mind cannot know the thing-in-itself, there was one ‘part’ of me which could; who I am. I myself am a thing-in-itself, and so if I can somehow penetrate my own experience, I can experience that which, to the subjective-objective mind, is forever out of reach. Kant had ruled this out; although he had speculated that I and the-thing-in-itself are the same , he, being foremost a rational thinker, denied that either were accessible. He never recognised that which is aware of thoughts and sensory intuitions; he reasoned us to the threshold of the unspeakable, but was unable to cross over. This led to one of the most basic flaws in Kant’s philosophy, first pointed out by a contemporary of Schopenhauer’s, Friedrich Adolf Trendelenburg who said that although Kant demonstrated very convincingly that the world as we know it is brought to us by mind, that doesn’t mean that the objective world is not still identical to it. We know, as Kant and Schopenhauer did, that there is something in the mind’s mind which does correspond to the ‘thing-in-itself’, but how do we know that reality doesn’t entirely correspond to it, that reality isn’t entirely mind-knowable? We don’t. We can put aside questions of coincidence (kind of astonishing that the mind should create a timespace reality in here that corresponded essentially to the one out there) because Trendelenburg’s objection can never be rationally addressed, or decisively closed for reasons that, as we shall see, Schopenhauer well understood, even if confusedly. Schopenhauer was able to cross over ‘Trendelenburg’s gap’ , but he was still dreadfully confused, a confusion which catastrophically coloured his output, leading him into dead-ends and absurdities. What saved him, grounded him, tuned him to frequencies of genius and opened up a new path for spiritual enquiry in the West was that, unlike Kant, he did not essentially begin his argument with mind — but with experience. He lived a far more interesting life than Kant and, crucially, was far more sensitive to that life. Schopenhauer himself tells us that what he writes is meaningful and interesting because it comes from direct experience, which then forms its conceptual system. Professional philosophy goes in the other direction, starting with concepts and systems and then arguing these into conclusions, which then contain nothing more meaningful or real than what is contained in their premises; which is to say nothing at all. This is why professional philosophy is dull, colourless and, in the end, meaningless, and it is why Schopenhauer is so instructive and entertaining, even when the divisions he makes between one field of experience and another are inaccurate, or even when he is just dead wrong, for it won’t be long before, even if expressed with misleading concepts or fitted into his flawed system, or coloured by his personal manias, a gleaming truth, direct from the source of life, radiates outwards. Because so much of what Schopenhauer says comes directly from life, it is instantly recognisable as true, it is very often funny and, wonder-of-wonders, it is clear. This, for Schopenhauer is, with the honorable exception of Kant, a mark of sincerity and therefore a sign that the writer is at least capable of divining the truth (for how can reality be discovered without sincerity?) rather than merely thinking. ‘…the genuine philosopher will generally seek lucidity and clarity and will always strive not to be like a turbid, raging, rain-swollen stream, but much more like a Swiss lake, which, in its peacefulness, combines great depth with a great clarity that just reveals its great depth. Clarity is the good faith of philosophers…’ (On the Fourfold Root of the Principle of Sufficient Reason and Other Writings) Schopenhauer rarely provides an idea without giving a concrete, and sometimes astonishingly vivid example. This is not just from a desire to be understood — although of course if someone has something meaningful to say, why would they not want to be understood? — but as a demonstration of the very point he is endeavouring to make, over and over again, that concepts come from life; or rather from a more direct sensate experience of life than thinking can provide. Conceptual thought without this experience is hollow, trivial, absurd and, no matter how well reasoned or argued, can never provide meaningful insight; as there was no such insight in the premises. ‘It all depends on where you have gotten your concept: if it is drawn from experience, fine, since its object exists and requires no further proof; in contrast, if it is hatched from your own half-brain, then all its predicates will not help it: it is just a figment of your imagination…’ (On the Fourfold Root of the Principle of Sufficient Reason) Schopenhauer (following Rousseau) was the first and greatest mind to fully tackle the principle fallacy of Western philosophy, the error that conceptual thought can grasp reality or, in its more extreme form, the outright lie, bordering on insane delusion, that conceptual thought somehow is reality. This madness, which has gripped Western minds since Plato, is as common today as it was when Schopenhauer was writing, even more so. Philosophers still mask their essential superficiality and emptiness with complex language (‘smug semantic trickery’), they still focus their attentions down ludicrously narrow avenues of enquiry, they still search for meaning in abstract reflection and they still get handsomely paid and widely praised for all this pointless activity. Schopenhauer would have been just as disgusted with the state of professional philosophy today as he was with that of his own day; with, as he said, those who live from philosophy rather than for it. Another hallmark of philosophical truth, although not one that Schopenhauer dwells on is, aside from its depth, its width. If a philosopher goes looking for reality where it can actually be found — everywhere — then he will necessarily make interesting discoveries everywhere. Not only will the essential truth of his investigations illumine whatever matter he directs his attention to, but that very light will glow from wherever his attention happens to be, whether that be on the body of his lover, a painting that is touching him, the river he sits next to, the feelings of anxiety he has on entering a room, the surprising boredom he feels after he has achieved what he wants, his relationship with his mother or a particularly funny friend he has — all of which, first of all, he will notice — for he is conscious — and, secondly, will naturally inform his philosophy, give him all kinds of things to contemplate and write about, metaphors in which to bridge distant fields of experience and greater insight into the single truth that animates all these disparate truths, all of which mark Schopenhauer’s work. Of course none of this matters in the slightest if the philosopher’s experience is limited by habit or by fear, if he lives his entire life in institutions, if he has no interest in matters of real existential importance — love and death — or if he is not truly sensitive; if everything that happens to him comes filtered through his emotions and his thoughts. To some extent this did apply to Schopenhauer; he was rich, not much of a lover, appallingly neurotic and although he travelled widely, not tremendously well-acquainted with the breadth of life. But as far as life is concerned, and acquaintance of it, he left all other thinkers far behind, before and since; with the notable exceptions of those who also had something to say (Hume, Nietzsche, Wittgenstein, etc.). Let’s return to Schopenhauer’s philosophy. Following Kant, he said that reality (or the thing-in-itself) comes to us via the senses, after which the mind, from that ‘raw’ sense-data, first of all creates perception — all the objects around us, existing in space and time, with their colours and shapes and so on — and then conception — abstract ideas about experience which are related in the mind into ideas, theories, memories, beliefs, systems and so on. First there is reality, ‘then’ there is percept, then there is concept. What this means is the Schopenhauer’s philosophy, like Kant’s, or like any philosophy which puts reality before perception and conception in this way, can never be proven. Discovering or establishing proof through rational means is a conceptual activity based on perception. That both percept and concept are secondary to reality is literally inconceivable, let alone unprovable. But if this is the case, how do we know that what anyone says about reality is truthful or accurate, and not complete nonsense? Schopenhauer gives something like the following analogy: Before the Rosetta stone was discovered it was impossible to understand Egyptian hieroglyphics. There was no ‘philosophy’ by which Egyptian symbols could be translated into languages which we understand. When the code was cracked, we had such a ‘philosophy’ which, almost at once, began to make sense of the hieroglyphics that we encountered. Does that mean we have now proved that our ‘philosophy of Egyptian’ is true? No, it just means that, up until now, it has worked. We can never prove that we understand Egyptian hieroglyphics — indeed that anyone understands any language. But we still know they do. We know that someone speaks fluent Chinese, for example, by their ability to deal with all the Chinese they encounter, just as we know that someone’s philosophy is true by its ability to deal with all the mysteries it encounters; not because it has been proved in logical thought. Schopenhauer does not speak well of logic and the quest for proof. ‘In general,’ he says ‘proofs are not for those who want to learn so much as for those who want to argue,’ and ‘learning logic for practical purposes is like training a beaver to build its own dam.’ Kant also had a low opinion of it: ‘Now one can take it as a certain and useful warning that general logic, considered as an organon, is always a logic of illusion, i.e., is dialectical. For since it teaches us nothing at all about the content of cognition, but only the formal conditions of agreement with the understanding, which are entirely indifferent with regard to the objects, the effrontery of using it as a tool (organon) for an expansion and extension of its information, or at least the pretension of so doing, comes down to nothing but idle chatter, asserting or impeaching whatever one wants with some plausibility.’ (Critique of Pure Reason) None of this means that seeking proof, or the rational activity that leads to it, is a useless activity. Schopenhauer is at pains to point out that for matters in which it makes sense to speak of cause and effect, time and space, abstraction and relation, the rational mind is the appropriate tool to use, and ignoring what it is and what it does in its proper sphere leads to confusion, and often to outright charlatanism. But thought is powerless to uncover that which precedes or encompasses conception and perception, which includes all the ‘ultimates’ upon which the rational mind is based. What happened before the universe was created? What actually is time and space? What is that which is conscious of the objects of the world? What is matter, the stuff of the universe? What is energy? What is that which my perceptions and conceptions are of? We can always approach the ‘why’ of things with the mind, never the what. Because thought (conception) is abstracted from perception, which itself is a ‘result’ of the thing-in-itself, thought can never reveal what experience really is, any more than a score can reveal what music really is. Thought can only tell us what it is, what thought is; abstract relation, conceptual symbol and conceptual law. Conception can tell us how all the ideas it presents to us are related with each other, it can form accurate ideas of cause and effect which consistently and usefully ‘map’ onto experience, it can, to some extent, preserve the past and predict the future, but it can never, ultimately, give us the slightest insight into that from which thought emerges, the terrain from which the map is taken. Schopenhauer provides a vivid (and funny) metaphor for this; ‘…given a complete [scientific] account of the whole of nature, the philosophical investigator would always feel like someone who finds himself in completely unfamiliar company without knowing how, and where each person in turn introduces another as his friend or cousin, thus making them sufficiently acquainted: in the meantime, however, the man himself, while expressing his pleasure at each new acquaintance, keeps asking himself: “But how the devil do I fit in to this whole group?”’ (The World as Will and Representation) What all this means is that there is no relation between reality and the experience of the self. Sensory experience in time and space can be said to ‘come from’ some thing-in-itself, but perceived knowledge cannot be related to the source of perceived knowledge because relation itself is a kind of perception-knowledge: ‘…the law of causality is the first condition of all empirical intuition, but this is the form in which all experience occurs: how then could the law first be derived from experience, for which it is the essential prerequisite?’ (On the Fourfold Root of the Principle of Sufficient Reason) This means that it is impossible to ever observe cause and effect. This was one of Hume’s most important ideas, taken up by Kant also. You can watch a million billiard balls hit a million more, but you’ll never see the ‘necessary connection’ which you think made it, along with everything else in the world that you observe, happen. Causality, insofar as we can speak or think of it, exists entirely in the mind. This means that there is never any logical reason to suppose that the next time you drop a rubber ball, for example, it won’t cease to bounce, or explode, or turn into a duck egg. The only reason you ‘know’ it won’t is because it’s bounciness ‘fits’ with all your other observations of the world, and with your innate time-space grasp of cause and effect. It is this interconnectedness of ideas which distinguishes waking reality from dreams. Dreams seem just as real to us as the waking world does, but when we wake we know it is a dream because nothing within the dream really ‘fits together’. Sometimes even in dreams we know this; we think ‘hang on, I don’t live with Woody Allen’, or ‘Why’s that narwal wearing a boa?’. This doesn’t mean that the waking world really is real though, just more real, or more reliable. Kant again: ‘The difference between truth and dream, however, is not decided through the quality of the representations that are referred to objects, for they are the same in both, but through their connection according to the rules that determine the connection of representations in the concept of an object, and how far they can or cannot stand together in one experience.’ (Prolegomena to Any Future Metaphysics) So what is really real? What is the ‘thing-in-itself’, and how do we experience it? Before we get to Schopenhauer’s highly misleading answer (and my clarification) it’s worth exploring a couple of foundational consequences of what he has said so far. The first is that if space and time really are functions of the mind, then there would appear to be, ultimately, or in some way, no difference between anything. There is no difference between the chair and the table, or between me and you, or even between you and Napoleon Bonaparte. ‘Difference’, in space (which we call position) and in time (which we call succession) do not really exist, and so everything is, ultimately, one. This doesn’t mean that there is no difference at all between you and Napoleon Bonaparte — probably you’ve noticed this? There must be something in the depths of youness which differentiates you from a dead emperor, but what that difference really is, the mind can never know, because it can only know the difference it perceives. What’s more there is no reason to rule out, along with this difference, a mind-blowing identity between you, at your core, and everything which is, has been, or ever could be. It is this identity which forms the essence of Schopenhauer’s ethics. He says that although we, you and I, are clearly somehow different, we are also, somehow, the same. Sensitivity to this sameness we call empathy, an extremely intimate but real — rather than merely reasoned — experience of the other. Because of the intimacy and reality of this experience — I feel, in hurting you, that I am hurting myself — it is both impossible to ignore and, ultimately, to define. All definable and abstract ‘ethical systems’ in contrast — including the demands of religions, the laws of states and all the ‘oughts’ that parents regularly assail their children with — are based on sand, and doomed. They need to be learnt with effort, enforced with effort and can effortlessly be circumvented. Another of the interesting consequences of understanding reality in the manner that Schopenhauer suggests is that the interminable debate about free will vanishes. This is because the world of percepts and concepts that we take as the world is the only world in which cause and effect can have any meaning. In the world — in the objective matter which I perceive and conceive — there is clearly no such thing as free will, or only in a limited sense. Everything ‘out there’ (and by ‘out there’ we also mean the world of the mind and of emotional motive) has a cause. We may, with difficulty, think it doesn’t, but we certainly do not live that way. When we get ill, when we see a magic trick, when someone tells us something bizarre, we immediately and quite sanely start looking for causes. We couldn’t live any other way. But this is not how it is when I look within, beyond my ideas and motivating emotions. Here I sense — and the feeling is unshakable — that I am essentially free, and from there I go on to assume that somehow everyone is, in their ‘core’, free also — and therefore responsible for what they say and do. I might philosophically assume that everyone’s acts are really determined — and of course there is some relative truth and value in that too — but such a defence doesn’t cut it when I am wronged or when someone is continually blaming the past for their misdeeds. Then we cannot help calling on their freedom. Living together is just inconceivable any other way. So whether I have ‘free will’ depends entirely on what aspect of ‘I’ am attending to. If it is the external I, the I that comes to me through percept and concept, then no, I am clearly not very free, if at all. But if I attend to that ‘from which’ percepts and concepts ‘arise’ — I discover an unshakable and, in the end, truthful sense of freedom. But how is this? And what is this ‘other kind’ of I? Earlier we saw, in Kant’s teaching, that what I perceive and conceive ‘comes from’ (note the scare quotes here and elsewhere in this account — it doesn’t really ‘come from’, but there’s no other way to talk about it) a thing-in-itself which I can never experience or rationally grasp, that is forever alien to me, remote. Kant however speculated, albeit fleetingly, that there is one thing-in-itself which we can experience: ‘…although extension, impenetrability, composition and motion — in short, everything our outer senses can transmit to us — are not thoughts, feelings, inclinations or decisions, and cannot contain them, as these are never objects of outer intuition, yet that same Something that grounds outer appearances and affects our sense so that it receives the representations of space, matter, shape, etc. — this Something, considered as noumenon (or better, as transcendental object) could also at the same time be the subject of thoughts…’ (Critique of Pure Reason) He abandoned this line of enquiry, but it was taken up enthusiastically (if misleadingly) by Schopenhauer who reasoned, correctly, that of course I am and must be a thing-in-itself too. Because space and time do not apply to the thing in itself, because distinction and separation are meaningless to it, because it doesn’t come into or pass out of existence, then I am, and have always been, that which I experience. It’s an insane idea of course, but logic leads towards it , instinct resonates with it, the finest minds in history and pre-history agree with it and it obliterates all the perplexing problems of philosophy. Take idealism and realism (also known, in another form, as the mind-body problem). Remember, idealism is the belief that nothing is real but the mind — that the ‘subjective’ I causes reality (or my experience of it, of the ‘body’), while realism (also known as materialism or, today, physicalism) is the belief that only the ‘objective’ world is real and that it, matter (or matter and energy), causes everything we perceive and experience, including ‘mind’. But if ‘cause’ cannot apply to the thing-in-itself, either the objective bodything out there I perceive or the subjective mindthing in here that I am, then both of these positions must be wrong. Cause and effect are created by the rational mind, so it is incoherent to say that objects or subjects cause the world. There is, ultimately, no such thing as cause. This has been reasonably clear since Kant suggested it — although in fact Kant made the foolish mistake of assuming that the thing-in-itself causes representation of it — and since Schopenhauer explained it, but philosophers (and in their own way, ordinary people) still cling to idealism (aka solipsism) and materialism/physicalism (or scientism). The reason for this, well understood by Schopenhauer and mercilessly exposed by Nietzsche, is that ultimately philosophy is not chosen for rational reasons, but for psychological ones. It is true that Kant has no answer to Trendelenburg, it is true that Schopenhauer offers no argument for his assumption that ‘all causality is only in the understanding and for the understanding,’ but these are not the reasons that the reality posited by Kant-Schopenhauer is rejected, any more than the absurdities of religion and science are the reasons why they are rejected. Philosophy — all philosophy, whether of the lectern, the pulpit or the pub — is psychology. Idealism is held to because excuses the self (nothing really matters because nothing is really real) and materialism/physicalism justifies the self (I’m not responsible, the objective world is). The reason that transcendental idealism has never been popular is that it sweeps away these two comforting excuses, which most people spend most of their lives alternating between. But what remains? What, if we remove the illusory self and the illusory world, is there? Here we are at the bedrock, at the thing-in-itself. What is it and how can we ‘know’ or experience it if not through percepts and concepts? And here we arrive at the fundamental flaw in the philosophical diamond that Schopenhauer gave the world. He calls the thing-in-itself will, a word which, at various times in his writings he identifies as the thing-in-itself (the timeless, spaceless ground of all existence), as emotion (inner sense; both conscious and unconscious), as evolutionary instinct (genetic / animal urges which manifest as emotions), as free will and as the subject. Five completely different things! This causes to him to make an enormous number of incoherent and contradictory statements. Christopher Janaway (who doesn’t really understand the heart of Schopenhauer’s philosophy) has highlighted many of the contradictions between Schopenhauer’s subject, will and the thing-in-itself, and there are many more. For example, Schopenhauer says that will (instinct) is a pitiless, selfish, violent, horrific and utterly blind force — a mad impulse to grow and feed and live, no matter what the cost, to produce living forms, or species, the sole purpose of which, if it can be called a ‘purpose,’ is to climb over other species in order to continue living and growing and spreading. The individual is meaningless to this ‘will’, it extinguishes them — us — without hesitation, in our millions. What’s more, living in us, it causes us untold misery, compelling us to steal, cheat, rape, kill, God knows what else, just to serve it. It is this incredibly bleak and bloodthirsty picture of reality which has earnt Schopenhauer the title of the ‘philosopher of pessimism’. Intelligently reading his books shows otherwise, but there is no question that he gets off on morbidly detailing this aspect of the will, and of the horrors of the world. But then, not long after, he says that the will can be ‘silenced’ or ‘negated’, that I can know the thing-in-itself in all its transcendent mind-obliterating totality. But what is it that does the silencing? The will? If so Schopenhauer is saying that will needs to silence will in order to know will. Sometimes he says that ‘reason’ or the ‘intellect’ must do the recognising and silencing — precisely that part of the self that he correctly identifies as the most superficial and illusory. Schopenhauer reveals his confusion most clearly when he talks about the ideas. These, which he picks up from Plato, are the eternal or archetypical essence of things which, when the will is silent, we see clearly in all their wonderful suchness. These ideas do seem to have some reality, although quite what that is isn’t clear; Plato’s entirely conceptual notion — that all actual horses are photocopies of a perfect abstract horse — is just loopy. More than likely the ideas are, as Schopenhauer suggests, a kind of intimation of eternity in the local forms we see around us; the quintessential woman that we love in our girlfriends perhaps, or the marvellous cowness of cow. But whatever they are — and perhaps we are just talking about the vivid, timeless reality of what is in front of our eyes, without the clouding distraction of self — what is it within us that experiences them? It can only be one’s own essence or consciousness. Again Schopenhauer considers essence to be a blind and miserable force, which, he says, must be denied in order for the intellect to perceive the ideas (or later, to perceive the eternal ground of existence), and so he is led into the same unintelligible philosophical loop; that will must be denied in order for will to perceive the will. Let’s clear it up then. We’ve got emotion, instinct, free-will and the ‘thing-in-itself’, all of which Schopenhauer calls ‘will’. We’ve also got percept and concept, which Schopenhauer understands very well, and consciousness, which sometimes he identifies with conceptualisation, or the subject, or as ‘the ground’ of consciousness. The truth is that consciousness is what I am, it is the thing-in-itself which I can experience, or be, and when I experience it, or am it, I somehow experience the ‘cause’ of reality everywhere and at all times. It’s not really a ‘cause’ — I experience it and its ‘relationship’ with the world as something which the mind can’t describe with its ordinary language. But in any case this thing-in-itself ‘enters’ reality as my self, as free-will, as feeling, as perception and as conception; as ‘my’ world. There once was a time when reality was lived in this way, an impossible consciousness manifesting as the self-world, but at some point the manifest self began to lead its own existence, pulling itself over my eyes, so to speak, and calling itself ‘me’. This false me (or ego), without original, unifying consciousness, was no longer a machine operated ‘by’ consciousness, but by itself; self-informed, or selfish. As the elements of self — will, feeling, perception and conception — lost their integration with consciousness, and the immeasurably nuanced context, they become crude and contentious. The subtle and harmoniously integrated I became a collection of degenerate parts; a limited range of selfish emotions, degraded perception, atrophied will and an endless stream of boring and predicable self-oriented thoughts: all the negativity which Schopenhauer mistakenly lumps together with the original I, the thing-in-myself, which he correctly identifies as the transcendent truth. Once we have clarified Schopenhauer’s thought in this way his errors become clear and pretty much everything else he says makes sense. When, for example, he declares that genius has a ‘a greater amount of cognitive power,’ what he means is that genius has a greater degree of consciousness combined with, necessarily, a more sensitive and, eventually, well-developed perceiving and thinking self. When he says that ‘genius is contrary to nature’ he means that consciousness is contrary to self-informed instinct. When he says that ‘the innermost kernel of the I’ persists after death while consciousness is an ‘epiphenomenon’ that passes away with death, he means that consciousness, the thing-in-myself, persists while the temporal self passes away. And when he declares that nature is a superviolent horrorshow while, at the same time, being the reflection of eternity, perfection incarnate, he means that nature perceived through the selfish self is constant warfare (just as matter is completely determined), while nature experienced directly, as the thing-in-myself, is paradise. There’s a lot more to say about Schopenhauer — his nonsense views on sex , women , the state and race , and, contrariwise his sublime observations on art, physiognomy and genius for example, or his remarkable love for and acute observation of animals, or his hilarious beef with the shallow confounder Hegel. I suggest you take a read yourself though. The World as Will and Representation is excellent, as is his large collection of essays, Parerga and Paralipomena (in both cases I recommend the recent Cambridge translations ). I’ll conclude part one of this investigation — we’ll return to how we might ‘know’ the thing-in-itself in part two — with one of Schopenhauer’s memorable images, this one expressing his conviction that death cannot possibly be the end of consciousness. As thing-in-itself it can no more die than it can be born: such notions are of the time-space world, of the time-space mind. Stepping back from that reveals a sense of eternity untroubled by beliefs, philosophies, gods and monsters. ‘Death is like the setting of the sun that only seems to be devoured by night, but in truth, as the source of all light, burns without pause, bringing new days to new worlds, forever rising and forever setting.’ (part 2 here) Notes

Assessment of young amblyopes. Array vs. single picture acuities. A linear array of intact and "jumbled" symbols suitable for testing young children has been used to evaluate 35 young amblyopes. Acuities with the jumbled array test were poorer than those obtained with single symbols: the discrepancy between the jumbled array and single picture acuities was greater for large acuity deficits than small deficits. The depth of amblyopia, indexed by the interocular difference in acuity, was greater for the jumbled array than single pictures. For the patients who could be tested with lines of Snellen letters the jumbled array acuities were closer to line letter acuities than single picture or grating acuities. These results suggest that the jumbled array provides a more sensitive test of amblyopia in young children than single symbols or gratings, possibly because the jumbled array presents contour interactions that exploit the crowding phenomenon.

Planet 9: If it exists, how did it end up at the edge of our Solar System? There is evidence to suggest there is a mystery planet lurking at the edge of our Solar System – but if it does exist, how did it end up there? Researchers from the Harvard-Smithsonian Center for Astrophysics (CfA) have studied several different scenarios for where Planet 9 came from, all of which are unlikely. Scientists have long speculated there is an unknown planet sitting at the edge of the Solar System. However, earlier this year researchers at the California Institute of Technology said they had found evidence to support this idea. They said the unusual orbital patterns of a group of Kuiper belt objects could best be explained by a giant planet gravitationally dominating that region of space. The unusual orbits of objects around Planet 9Caltech/R. Hurt IPAC With mounting evidence to suggest Planet 9 does indeed exist, scientists have put forward four scenarios to explain where it came from. "The evidence points to Planet 9 existing, but we can't explain for certain how it was produced," says CfA astronomer Gongjie Li, lead author on one of the three studies submitted to the Astrophysical Journal Letters. Planet 9 orbits the Sun at between 40 billion to 140 billion miles. This is far beyond any of the other planets in the Solar System. To explain its position, scientists ran computer simulations to look at three ways the planet ended up where it has. In the first, Planet 9 was pulled outwards by a passing star, which had a stronger gravitational pull on it than the Sun. This would have tugged the planet into a wider and more elliptical orbit. As the Sun was born in a cluster of stars, the likelihood of another star passing nearby in the Solar System's early history is the most plausible explanation. However, the simulations revealed the chance of this happening is only 10%, as a star would probably pull Planet 9 out of the Solar System completely. Researchers find a possible ninth planet beyond NeptuneIBTimes UK In another simulation, run by CfA's Scott Kenyon, scientists said Planet 9 may have formed much closer to the Sun. After interactions with Jupiter and Saturn, the planet was effectively kicked out to its current position. A third scenario involves Planet 9 forming at the far reaches of the solar system in the first place. The team said under the right conditions, the hypothetical passing star (mentioned earlier) would be enough for it to be pushed into its current orbit. A final, and very unlikely scenario, involves Planet 9 being an exoplanet that either pulled into the orbit by a passing star, or drifted into the orbit by accident. This likelihood of this scenario is 2%, however. Whether or not the planet even exists is still in question, however. As Jim Green, director of Nasa's Planetary Science Division, said earlier this year: "The possibility of a new planet is certainly an exciting one for me as a planetary scientist and for all of us. This is not, however, the detection or discovery of a new planet. It's too early to say with certainty there's a so-called Planet X. What we're seeing is an early prediction based on modelling from limited observations. It's the start of a process that could lead to an exciting result."

My group and I are carrying out an investigation to find out how different concentrations of glucose affect osmosis in potato cells. Osmosis is the passive process of diffusing water which means within this experiment particles of water will move from a lower concentrated area to an area of higher concentration ; in this case it will be a concentration of glucose. Diffusion is when particles move from an area of high concentration to low concentration ; for example when a smell of cooking spreads through the house. Osmosis is a vital mechanism in the transport of fluids in living organisms. This point is clearly proven when you put a plant cell in water.Osmosis causes water to enter the plant cell, which results in swelling. However, the cell will not burst as the plant cell walls are made up of an extremely strong substance called cellulose. The swelling within the cell eventually stops, and at this point, the cell is said to be turgid. This process involving osmosis is important within a plant as it allows the stems to become strong and upright. There are many other cases where osmosis is responsible for the basic survival of living organisms. Here is a diagram which visually shows how osmosis will occur in this situation : In my group we are going to perform our investigation by firstly taking five equally sized test tubes and a potato which we will take a sample of using a cork borer. Afterwards, we will cut five equal pieces of potato using a knife, so that they are more or less the same weight. Once we have weighed all the five pieces of potato, we will get the five different glucose solutions that will be used throughout the experiment. These solutions will consist of different molars of glucose. For example, there will be a solution of 0 molars, 0.2 molars, 0.4 molars, 0.6 molars and 0.8 molars. Using a measuring cylinder, we will measure accurately 10 ml of each solution and put it in the appropriate test tube which would have been labelled relavant... YOU MAY ALSO FIND THESE DOCUMENTS HELPFUL ...INTRODUCTION Osmosis is the movement of water molecules from high concentration to low concentration through semipermeable membranes, caused by the difference in concentrations on the two sides of a membrane (Rbowen, L.). It occurs in both animals and plants cells. In human bodies, the process of osmosis is primarily found in the kidneys, in the glomerulus. In plants, osmosis is carried out... ... Aim To investigate the effects of increasing salinity on potatocell mass. Background Information This experiment is based upon osmosis. Osmosis can be defined as the net movement of water molecules from a region with high concentration to a region with low concentration. This movement must take place across a partially permeable membrane such as a... ... Introduction In this experiment I am going to investigate the effect of varying concentration of a differing glucose solution on the amount of osmotic activity, between the solution and a potato tuber of a given size. The purpose of this experiment is to demonstrate how living cells rely on osmosis, the diffusion of water. ... ...AN INVESTIGATION INTO THE EFFECTS OF OSMOSISPOTATOCELLS SKILL AREA P: PLANNING Introduction The aim of this investigation is to see the effect of varying concentrations of glucose solution on the amount of osmotic activity between the solution and a potato chip. An investigation into Plasmolysis in onion cells was undertaken prior to this... ...EXPERIMENT TO SHOW THE EFFECT OF DIFFERENTCONCENTRATION OF GLUCOSE ON POTATO STRIPS INTRODUCTION: Molecules of liquid and gas are constantly in motion, they move randomly in all directions and bounce around in all directions and bounce around and into each other. As they move, they tend to spread out moving from areas with many molecules to areas with fewer molecules . This process of spreading out is called... ...glass was filled with cold tap water. A drop of red food coloring was dropped in the cup. A stopwatch was used to measure the time it took for the food coloring to get to the bottom of the cup. The average diffusion rate was .78 cm a second. If a different color was used, I do not think it would have made a difference in the results. The mood of the person experimenting could possibly alter the results. For example, if a person is under stress, they may accidentally squeeze too... ...I am going to investigate osmosis when potato is placed in differentconcentrations of sucrose. I am aiming to witness osmosis in 5 differentconcentrations of sucrose. I will use 5 varying concentrations so that I have a wider spread to compare the results, and check that I don’t have any anomalies Prediction Osmosis is the process of diffusion of...

Currently, the wireless communication systems such as the mobile telephone systems or the wireless LANs (Local Area Networks) are widely used. In order to further improve the transmission rate and the transmission quality in the wireless communication technologies, next-generation technologies are being actively discussed. Some wireless communication technologies enable communication between two wireless communication devices (e.g., between a wireless base station and a mobile station) by use of multiple carriers. In some cases, the use of multiple carriers is called carrier aggregation, and each carrier in the carrier aggregation is called a component carrier. In wireless communication, control signals are transmitted from a first wireless communication device to a second wireless communication device. The information transmitted by the control signals may include information which is to be referred to by the second wireless communication device for receiving data transmitted from the first wireless communication device (e.g., information indicating the format used in the data transmission), or information designating a transmission procedure to be used by the second wireless communication device in transmission of data (e.g., information designating a format to be used in the data transmission). For example, information transmitted from a wireless base station to a mobile station may include information to be referred to by the mobile station when the mobile station receives data through a downlink data channel, or information to be referred to by the mobile station when the mobile station transmits data through an uplink data channel. In the case where wireless communication is performed between two wireless communication devices by use of multiple carriers, a wireless resource used in transmission of a control signal may belong to a carrier different from the carrier to which a wireless resource used in data transmission controlled by the control signal belongs. In this case, the problem is the method for recognition of the carrier to which the control signal is applied. According to a technique which has been proposed as a recognition method, the carrier to which the control signal is applied is clearly indicated by attaching a bit for carrier recognition (carrier indicator) to the control signal. (See, for example, “Final Report of 3GPP TSG-RAN WG1 #57 v1.0.0”, 3GPP (3rd Generation Partnership Project) TSG RAN WG1 #57bis, R1-092292, July 2009, Section 15.4.) Another recognition method which has been proposed is applied to a wireless base station scrambling a CRC (Cyclic Redundancy Check) bits attached to a control signal, with a scrambling sequence corresponding to an ID (Identifier) assigned to a mobile station as an opposite party. (See, for example, “Control signaling for carrier aggregation”, 3GPP TSG RAN WG1 #55bis, R1-090375, January 2009.) According to the above recognition method, IDs in the number of carriers are assigned to each mobile station, and the CRC bits are scrambled with a scrambling sequence corresponding to one of the IDs. Therefore, it is possible to concurrently recognize both of the mobile station and the carrier to which the control signal is applied. However, according to the technique disclosed in 3GPP TSG RAN WG1 #57bis, R1-092292, the wireless resources are excessively spent for attaching the carrier indicator to the control signal, so that the transmission efficiency of the control signal is lowered. In addition, according to the technique disclosed in 3GPP TSG RAN WG1 #55bis, R1-090375, the IDs in the number of carriers are assigned to the opposite party. Therefore, the ID exhaustion becomes a problem. In the case where the length of the bit sequence of the ID is fixed, the number of opposite parties with which communication can be performed in parallel is limited. For example, in the case where five carriers are available in total, the upper limit of the number of opposite parties with which communication can be performed in parallel can become one fifth of the number of opposite parties with which communication can be performed in parallel when only one carrier is used in communication with each opposite party. Further, in some cases, information indicating a border between the control signal and a data channel is transmitted as a separate control signal. In such cases, when the separate control signal is not correctly received, the border between the control signal and the data channel cannot be correctly detected, so that data cannot be correctly received through the data channel.

Dipole antennas have long been used in various communications systems, including radio, television, and radiotelephone systems. It is well known that the lengths of the dipole arms on the antenna should be adapted to the wavelengths (.lambda.) of the signals transmitted and received. Typically, a plurality of arms having different lengths are used, in order to cover a predetermined range of frequencies. The sequence and spacing of these arms, and of any reflector behind them, determines various characteristics of the resulting beam or radiation field. These characteristics include vertical beam width, horizontal beam width, and front-to-back (F/B) ratio, i.e. the ratio of signal strength in front of the antenna to signal strength in back of the antenna. When a number of different arms are used, each arm makes its own contribution to the resulting field, and the overall expected result rapidly becomes difficult to calculate mathematically in advance. Therefore, considerable experimentation is often needed to achieve desired beam characteristics. A well-known log periodic dipole antenna (LPDA) design is the "tree" configuration, in which parallel arms extend sideways from a central "trunk" or "standoff," the bottom arm near the base is the longest, and each successive arm is shorter toward the top of the antenna. Such LPDA designs typically result in a front-to-back (F/B) ratio less than 40 dB. This F/B ratio is considered insufficient for use in current PCS (Personal Communication System) cellular telephone sites, since radiation emanating out the back of the antenna tends to cause interference among adjacent sites. A horizontal beam width of 90 degrees is typical. However, in highly congested urban environments, it is preferable to have horizontal beamwidth of 65 degrees, which is obtained by using two parallel columns of dipoles, spaced 0.25 .lambda. to 0.30 .lambda. apart. The wavelength lambda (.lambda.) is the inverse of the frequency. The frequency band allotted for PCS use in the United States is between 1.85 GigaHertz and 1.99 GigaHertz, with a center frequency 1.92 GHz. The PCS band allotted in Europe has a center frequency 1.78 GHz, meaning that the wavelength is about 8% greater. Accordingly, antenna dimension examples stated for the U.S. should be scaled up about 8% for use in Europe. My earlier LPDA design work has included an "hourglass" dipole strip configuration, in which top and bottom arms are longer than one or more middle arms. This design works well for generating a 90 degree beamwidth, but when used for generating a 65 degree beamwidth, typically results in F/B ratios in the range between 37 dB and 42 dB, better than provided by the "tree" configuration, but still insufficient. Reference is made to pending application U.S. Ser. No. 08/807,560 by myself and a colleague.

Domex Philippines joined World Toilet Day on November 13 According to a study conducted by the World Health Organization (WHO), children around the world miss an estimated 443 million school days each year because of diseases caused by poor sanitation and hygiene practices, including using unsanitary toilets. Simply put, this number of sick days is equivalent to all grade school and high school classrooms in the Philippines being empty for one month. This sanitation-related concern hinders the children’s learning and significantly reduces their quality of life. Illnesses due to poor toilet sanitation include diarrhoeal disease and parasitic worm infections which can lead to nutritional deficiencies, physical and mental stunting, and death. Every day, millions of children in rural communities and urban households nationwide are exposed to these health problems by using dirty toilets populated by disease-carrying germs. According to Dr. Luisa Efren of the Philippine Public Health Association (PPHA), said that the simple act of proper toilet sanitation can help prevent the spread of these germs.Children continue to be at risk from these deadly diseases because many households still use ordinary laundry bleach to clean their toilet bowls, which is not enough to kill all toilet germs. “A toilet that looks clean to the naked eye may not necessarily be free from germs and bacteria. So, it is very important to properly sanitize all surfaces using a germ-kill expert with proven efficacy in eliminating bacteria, and not just any ordinary laundry bleach,” Dr. Efren warned. For the past three years, Unilever Philippines through its germ-kill expert brand Domex, has been working with UNICEF and PPHA to champion the One Million Clean Toilets Movement. This advocacy program aims to educate Filipinos on proper toilet hygiene and the need for sanitized household toilets to keep their families safe against disease-causing germs. On November 13, World Toilet Day, join the One Million Clean Toilets Movement in the fight against disease causing germs to save lives, one clean toilet at a time. To learn more about Domex and how you can stay safe from disease with a germ-free home, visit Domex Philippines on Facebook.

Desperate parents of nail-biting kids and furniture-chewing dogs often turn to foul-tasting compounds to deter their loved ones’ gnawing. In the future, we might be trying something similar with mosquitoes, as scientists now say that malaria-causing mosquitoes use scent and taste to decide who to bite. They published their research in the journal Nature Communications. The mosquito, that scourge of summer, is more than just a nuisance. Mosquito-borne viruses like Zika and dengue are on the rise. Rates of malaria are down, but still high; around 214 million people were affected in 2015 alone. As drug developers race to develop vaccines, other scientists are hoping to find ways to keep disease-carrying mosquitoes from biting in the first place. To do that, they’ve got to know mosquitoes inside and out. A mosquito lives by its noses. Yes, noses, plural. Every mosquito has three sets of scent-detecting parts: two antennae, two fuzzy mouthparts called maxillary palps, and two little spots called labella at the end of its proboscis. The antennae and palps are scent-only, but the labella contain neurons for sensing both smell and taste. That’s a lot of olfactory information for a teeny tiny brain to take in. To find out how the mosquitoes do it, researchers tinkered with the genes of the malaria-carrying mosquito Anopheles gambiae. They tricked out the mosquitoes with a gene that would cause cells called odorant receptors (ORs) to glow bright green, which would make them easier to spot under a microscope. Building fluorescent proteins into bug parts is not a new technique, but it’s never been done before in mosquitoes. This is a female Anopheles gambiae mosquito with olfactory neurons on the antennae, maxillary palp, and labella labeled in green. Image Credit: Olena Riabinina and Courtney Akitake, Johns Hopkins Medicine By looking at the mosquitoes’ glowing ORs, the team was able to trace the paths from the pests’ sense organs to their brains. They found that information taken in by the antennae and maxillary palps was sent to brain regions called antennal lobes (this process is the same in flies). But signals from the labella went over to an area called the subesophageal zone—an area that had previously only been associated with taste. The researchers say this likely means that a mosquito not only sniffs us but tastes us, too, poking with the end of its proboscis to confirm we’re edible before unsheathing its gross, syringe-like feeding needles. It’s an unsettling concept, to be sure, but it might just help us out down the road. Co-author Christopher Potter, a neuroscientist at Johns Hopkins University, says we could use An. gambiae's brain cells against it—all we have to do is convince it that we taste terrible. “Our goal is to let the mosquitoes tell us what smells they find repulsive and use those to keep them from biting us,” he said in a statement. Lead author Olena Riabinina, now at Imperial College London, noted that their success with the glowing protein has created new possibilities for mosquito research. “We were pleasantly surprised by how well our genetic technique worked and how easy it is now to see the smell-detecting neurons,” she said. “The ease of identification will definitely simplify our task of studying these neurons in the future." Know of something you think we should cover? Email us at tips@mentalfloss.com.

Folate levels in human liver from autopsies in Canada. Liver folate is considered to be a direct measure of body folate stores. Specimens of 560 livers were therefore collected during autopsies at hospitals in seven cities across Canada, and analyzed for "total" folate. The values obtained ranged from 2.7 to 15.6 micrograms/g. Stillborn infants had the lowest mean liver folate levels (5.9 microgram/g). Mean liver folate levels tended to increase from birth to a peak between 11 to 30 years (8.8 micrograms/g), and then decreased with increasing age. This trend was reflected in an increased proportion (6 to 15%) of folate levels in the 3.1 to 5.0 microgram/g range in older subjects. Liver folate levels of male subjects did not differ from those of females. Mean liver levels tended to be slightly higher, and the number of liver levels between 3.1 to 5.0 micrograms/g was lower, in subjects dying accidentally (7.8%) and from heart and coronary disease (8.7%), than in those dying from cancer, respiratory diseases or other causes (11.7 to 15.3%). Only two subjects had liver folate levels below 3 microgram/g. In this survey, there was thus little evidence of severe folate deficiency.

Battle of Portland Harbor The Battle of Portland Harbor was an incident during the American Civil War, in June 1863, in the waters off Portland, Maine. Two civilian ships engaged two vessels under Confederate States Navy employment. Background Around June 24, a Confederate raider named the Tacony, commanded by Lieutenant Charles Read, CSN, was being pursued by the Union Navy. To thwart their pursuers, at about 2 AM on June 25, the Confederates captured the Archer, a Maine fishing schooner out of Southport. After transferring their supplies and cargo onto Archer, the Confederates set fire to Tacony, hoping the Union Navy would believe the ship was destroyed. On June 26, a Confederate raiding party entered the harbor at Portland late in the evening, sailing past Portland Head Light. The rebels disguised themselves as fishermen; they planned to try to destroy the area's commercial shipping capability, and then to escape out of the harbor. Battle When the raiders left the port area on June 27, they proceeded to the federal wharf. Having the advantage of surprise, the crew seized a cutter belonging to the Revenue Service, the USRC Caleb Cushing (named for a Massachusetts congressman, United States Attorney General and Minister to Spain). Their original intent was to seize the side wheel steamer Chesapeake, but its boilers were cold. As they would lose too much time in getting the steam up, they took Cushing. They escaped and sailed out to sea. News spread of the Confederate actions and the Army garrison at Fort Preble in nearby South Portland was alerted to the rebel intrusion. The Confederates had been observed by several persons while taking over the cutter, and public fury was aroused. Thirty soldiers from Fort Preble were assigned to pursue the raiders; they took a six-pound field piece and a 12-pound howitzer. Accompanied by about 100 civilian volunteers, the soldiers commandeered the steamer Forest City, a side-wheel excursion ship, and the Chesapeake, whose steam was finally up. All of the civilians on board were issued muskets to defend against the Confederates. Forest City, the faster ship, was the first to catch up to Cushing and Archer. Cushing opened fire on Forest City when it had come within the range. The captain of Forest City was afraid to pursue any further. Cushing, being a revenue cutter, had two secret compartments hidden in the captain's stateroom. Confederate Lieutenant Read had not discovered the cache of powder and ammunition stored there. If he had, the outcome could have been very different. Chesapeake, which had left port sometime after Forest City with Portland's Mayor Jacob McLellan in command, finally caught up and continued on toward Cushing. The wind was beginning to blow against the Confederate sailors and the steamers soon caught sight of Cushing. Read, the Confederate lieutenant, ordered Cushing torched; its munitions exploded after the ship was abandoned by her twenty-four crewmen, who escaped in lifeboats. They surrendered to Mayor McLellan and were held as prisoners of war at Fort Preble. Archer was also soon captured, and all the rebels were returned to Portland. Aftermath It was discovered that the Confederates were in possession of over $100,000 in bonds. These were to be paid after a treaty for peace was ratified between the North and the South. Public anger against the Southerners was high, and the city requested additional troops to safeguard the prisoners. When they were to be transported to Boston in July, the men had to be spirited out of Portland during the night to prevent a riot from breaking out. They were removed to Boston Harbor, and held at Fort Warren. Footnotes References Maine Bureau of Corporations, Elections, and Commissions Harper's Weekly, 11 July 1863 Confederate Navy Research Center, Mobile, Alabama The New York Times, 28 June 1863. Portland Harbor Portland Harbor Portland Harbor Portland Portland Harbor Category:19th century in Portland, Maine Category:June 1863 events

Sample Surveys and the 1940 Census After a 72-year wait required by law, the National Archives has released individual records from the 1940 Census, opening a gold mine for people researching their family histories. But the 1940 Census also played a notable role in the history of census-taking: It helped usher in the modern era of sample surveys. With the nation deep in the Great Depression, government officials planning for the 1940 Census hoped to expand the topics they asked about in order to guide federal policy-making in an era of expanded government. They wanted to add questions to the census about people’s incomes, migration histories and housing conditions. But adding these new questions would have made the census form too lengthy, and been too costly. After lobbying by some of its younger statisticians, the Census Bureau’s solution was to turn to the new science of survey sampling. To make room for new topics, the Census Bureau decided for the first time to ask some questions of a sample of respondents, not of everyone. Those responses were extrapolated to the total population. New Questions The 1940 Census included 34 population questions and 31 housing questions for the general population, and an additional 16 questions (on six topics) that were asked of a sample of respondents. Some questions asked of the sample had been asked in previous counts, and others were asked for the first time in 1940. The housing questions were listed on a separate form, but in practice enumerators asked them along with the population questions. Census Bureau officials knew that some of the new questions would be controversial, and that some people would find them intrusive. There has been debate since the first census, in 1790, about how much it is appropriate for the government to ask beyond a basic count. In 1940, not many Americans had previously been asked to give financial information to the federal government, according to historian Margo J. Anderson’s book, “The American Census: A Social History.” Only a minority of Americans then filed income tax forms; only 15 million forms were filed in 1940, when the national population was counted at 132 million. Census officials knew that the strongest opposition to the income question came from high-income Americans, but the information they needed most was from low- and middle-income Americans. So census-takers were told to ask for actual wage or salary amounts only up to $5,000 a year; at the time, three-quarters of American families made no more than that. Incomes of $5,000 or more would be reported as “$5,000+.” The instructions to census-takers said: “Some persons who might otherwise be reluctant to report wages or salary would be quite willing to do so if they learn that the amount above $5,000 need not be specified.” If people were reluctant to give their income information to the enumerator who knocked on their door, they could choose to fill out a card that would be sealed in an envelope and mailed to the Census Bureau. The question on income was toward the end of the questionnaire, “because if the enumerator got kicked out of a household when that question was asked, the interviewer would have already obtained the answers to the previous questions,” recalled Census Bureau official Edwin D. Goldfield, in an interview for a bureau oral history archive. The housing questions amounted to the first national inventory of housing, according to Anderson. They included questions about the home’s water supply and toilet, if any. There also was a series of questions about the home’s value and cost. The housing forms were destroyed, and are not part of the 1940 Census release of records. Sample Survey Since the first U.S. census, the same set of questions had been asked of everyone. But in 1940, the desire for additional data coincided with improvements in survey methodology and theory that allowed the Census Bureau to add more questions without burdening the entire population. Statisticians in the 1930s had made great advances in designing and implementing sample surveys, in which a randomly selected subset of respondents supply data that are extrapolated to represent views of the entire population. The importance of drawing a random sample was made clear after Literary Digest, a well-known magazine, conducted a straw poll of respondents whose names were obtained mainly from automobile registration lists and telephone books. The poll, which had a gigantic response of two million postcards (about 1,000 is considered an adequate sample today), indicated a landslide victory for Republican Alf Landon over Democrat Franklin D. Roosevelt in the 1936 presidential election. Although factors other than the flawed sample also played a role in that failed forecast, the Literary Digest debacle was a strong force in the rise of surveys based on scientific probability samples, in which any American adult has a known chance of being asked for a response. Other polls taken in 1936 that were based on random samples were much more accurate than the Literary Digest poll. Some long-time Census Bureau officials had resisted incorporating scientific sampling into the decennial census, believing it would “downgrade the validity of census information, because you had to say that this is based on a sample,” recalled Ross Eckler, a former Census Bureau director who was interviewed for a Census Bureau oral history project. Armed with information from some smaller government sample surveys, a younger generation of statisticians persuaded their skeptical elders at the Census Bureau that such an approach should be part of the decennial census. The introduction of sampling did not cause any notable public controversy or challenge from members of Congress, according to Eckler and other sources. In the 1940 Census, the sampled population consisted of every 20th person interviewed by any given enumerator. In practice, enumerators filled out forms that had space for information from 40 people on each side, and were told to ask the supplemental questions of people whose names fell on certain designated lines. According to the instructions for enumerators, the questions should be asked about anyone whose name was listed on a certain line, “whether this be the head, his wife, a son or daughter, an infant, a lodger, or any other member of the household.” The sample questions included place of birth of the person’s mother and father; “mother tongue” in the household during early childhood; three questions about military service and veteran or veteran-family status; three questions about Social Security receipt; and questions about occupation, industry and class of worker. Women who were or had been married were asked whether they had been married more than once, how old they were when first married and the number of children they had ever had. Sampling Since 1940 Every decennial census from 1950 to 2000 included widened use of sampling to ask additional questions of part of the population, which eventually were asked on a separate form known as the “long form.” Since 2000, the American Community Survey, a sample survey, has asked the same questions on a continuous basis that the census had asked once a decade on the long form. The Census Bureau also uses sample surveys to evaluate the quality of the census and to test new questions and survey methods. The techniques for drawing a sample, and extrapolating results, have become increasingly complex. Over the years, the Census Bureau moved from a person-based sample to a household-based sample. Instead of being drawn from people interviewed by one enumerator, samples were drawn from lists of addresses. The reliability of sample surveys is now taken for granted within the Census Bureau, and wider research community, where sample surveys are the basis for measuring public opinion, political attitudes, employment levels and other information. Some uses still can generate controversy, however. Since 1950, the Census Bureau has taken a post-enumeration sample survey to check the quality of the full count. In 2000, the Census Bureau planned to use a sample survey as the basis not only for checking the accuracy of its original count, but also for amending that official count if it was found that the post-census survey data would improve the quality of the enumeration. The agency planned to conduct a post-enumeration survey, match the results with actual census records and apply a statistical technique known as dual-systems estimation to correct flaws in the original enumeration. However, the application of survey sampling to produce the official counts used to apportion congressional seats among the states met opposition, mainly from Republicans who expressed concern that results could be manipulated for partisan purposes. The use of survey sampling for apportionment purposes was successfully challenged in the U.S. Supreme Court before the 2000 Census was taken. The Census Bureau did take a post-enumeration survey in 2000, and considered trying to use its results for other purposes, such as producing data for redistricting within states or allocating federal funds among states and localities. But problems were discovered with the survey results, so they were not used for other purposes. The bureau did conduct a post-enumeration survey after the 2010 Census, called the Census Coverage Measurement program. It is intended to produce measurements of undercount, duplicate count and other error, but not to amend the count. Results are expected sometime this year. About Pew Research Center Pew Research Center is a nonpartisan fact tank that informs the public about the issues, attitudes and trends shaping America and the world. It conducts public opinion polling, demographic research, media content analysis and other empirical social science research. Pew Research Center does not take policy positions. It is a subsidiary of The Pew Charitable Trusts.

Q: What does "use strict" do in JavaScript, and what is the reasoning behind it? Recently, I ran some of my JavaScript code through Crockford's JSLint, and it gave the following error: Problem at line 1 character 1: Missing "use strict" statement. Doing some searching, I realized that some people add "use strict"; into their JavaScript code. Once I added the statement, the error stopped appearing. Unfortunately, Google did not reveal much of the history behind this string statement. Certainly it must have something to do with how the JavaScript is interpreted by the browser, but I have no idea what the effect would be. So what is "use strict"; all about, what does it imply, and is it still relevant? Do any of the current browsers respond to the "use strict"; string or is it for future use? A: This article about Javascript Strict Mode might interest you: John Resig - ECMAScript 5 Strict Mode, JSON, and More To quote some interesting parts: Strict Mode is a new feature in ECMAScript 5 that allows you to place a program, or a function, in a "strict" operating context. This strict context prevents certain actions from being taken and throws more exceptions. And: Strict mode helps out in a couple ways: It catches some common coding bloopers, throwing exceptions. It prevents, or throws errors, when relatively "unsafe" actions are taken (such as gaining access to the global object). It disables features that are confusing or poorly thought out. Also note you can apply "strict mode" to the whole file... Or you can use it only for a specific function (still quoting from John Resig's article): // Non-strict code... (function(){ "use strict"; // Define your library strictly... })(); // Non-strict code... Which might be helpful if you have to mix old and new code ;-) So, I suppose it's a bit like the "use strict" you can use in Perl (hence the name?): it helps you make fewer errors, by detecting more things that could lead to breakages. Strict mode is now supported by all major browsers. Inside native ECMAScript modules (with import and export statements) and ES6 classes, strict mode is always enabled and cannot be disabled. A: It's a new feature of ECMAScript 5. John Resig wrote up a nice summary of it. It's just a string you put in your JavaScript files (either at the top of your file or inside of a function) that looks like this: "use strict"; Putting it in your code now shouldn't cause any problems with current browsers as it's just a string. It may cause problems with your code in the future if your code violates the pragma. For instance, if you currently have foo = "bar" without defining foo first, your code will start failing...which is a good thing in my opinion. A: The statement "use strict"; instructs the browser to use the Strict mode, which is a reduced and safer feature set of JavaScript. List of features (non-exhaustive) Disallows global variables. (Catches missing var declarations and typos in variable names) Silent failing assignments will throw error in strict mode (assigning NaN = 5;) Attempts to delete undeletable properties will throw (delete Object.prototype) Requires all property names in an object literal to be unique (var x = {x1: "1", x1: "2"}) Function parameter names must be unique (function sum (x, x) {...}) Forbids octal syntax (var x = 023; some devs assume wrongly that a preceding zero does nothing to change the number.) Forbids the with keyword eval in strict mode does not introduce new variables Forbids deleting plain names (delete x;) Forbids binding or assignment of the names eval and arguments in any form Strict mode does not alias properties of the arguments object with the formal parameters. (i.e. in function sum (a,b) { return arguments[0] + b;} This works because arguments[0] is bound to a and so on. ) arguments.callee is not supported [Ref: Strict mode, Mozilla Developer Network]

The "normal" adult human has 32 teeth. An adult should have 32 teeth with 8 incisors, 4 canines, 8 premolars, and 12 molars which include the four wisdom teeth. Humans develope 2 sets of teeth in their lives. the first set known as baby teeth began to come though at about age 6 months. They eventually have about 20 teeth as children before they began to lose ... Doctors say to wait 12 days after a missed period to take a test but if you no your regular with your period then you can take a test a day after you miss your period. The best time to confirm pregnancy is 3 to 4 days after you miss your regular period. If you test too early your pregnancy test may show negative result even if you are pregnant, therefore wai ... There are 14 systems that make up the human body. The Circulatory or Cardiovascular System (Heart, Blood Vessels (Arteries, Capillaries, Veins).) The Digestive or Gastrointestinal System (Mouth, Esophagus, Stomach, Small intestine, Large intestine) The Endocrine or Glandular or Hormonal System (Hypothalamus, Pituitary, Thyroid, Parathyroids, Adrenals, Pineal ... Protein is one of the basic building blocks of the human body, making up about 16 percent of our total body weight. Muscle, hair, skin, and connective tissue are mainly made up of protein. However, protein plays a major role in all of the cells and most of the fluids in our bodies. In addition, many of our bodies' important chemicals enzymes, hormones, neuro ... Approximately 21.5 million American adults, or about 9.5 percent of the U.S. population age 18 and older in a given year, have a depressive disorder. Nearly twice as many women (12.0 percent) as men (6.6 percent) are affected by a depressive disorder each year. These figures translate to 13.7 million women and 7.8 million men in the U.S. Pre-schoolers repres ...

1. Technical Field The inventions relate to optical communication systems. 2. Discussion of the Related Art This section introduces aspects that may be helpful to facilitating a better understanding of the inventions. Accordingly, the statements of this section are to be read in this light and are not to be understood as admissions about what is in the prior art or what is not in the prior art. Some current communications systems propose the use of pulse position modulation (PPM) to achieve high sensitivity at a receiver. In PPM, a signal is transmitted during a single time slot of the M time slots within each symbol period. In each symbol period, the receiver examines each of the M time slots to identify the time slot in which the signal wits transmitted. Unfortunately, increasing the number of time slots per symbol interval causes the PPM pulses to become temporally narrow and to have correspondingly larger bandwidths. For that reason, PPM may not provide a high spectral efficiency.

International tools To promote and protect human rights, States draw up and adopt standards on the subject within different international and regional forums. Human rights agreements impose obligations on the States which ratify them. They also often establish control mechanisms for the observance of these obligations. The Universal Declaration of Human Rights The Universal Declaration of Human Rights was adopted in 1948 by the General Assembly of the UN. For the first time, the different human rights were listed and proclaimed for all humankind. Since its adoption, this Declaration has been a referential text of great moral value. The great principles and standards established in this basic text have since been codified in a series of binding international and regional instruments. These have enabled us to develop and specify human rights standards. International and regional agreements Since the adoption of the Universal Declaration of human rights, several specific agreements on fundamental rights or groups of human beings have been drawn up. Since then, most of these agreements have been followed by optional protocols which set out the States’ obligations on each of these issues even more specifically. The standards drawn up at UN and Council of Europe level are extremely close and complementary in content. At UN level: Belgium is party to the main human rights agreements and optional protocols. Below is an overview of these different instruments: Agreements ratified by Belgium: International Covenant on Civil and Political rights (ratified on 21 April 1983) International Covenant on Economic, Social and Cultural rights (ratified on 21 April 1983) International Convention on the Elimination of All Forms of Racial Discrimination (ratified on 7 April 1975) Convention on the Elimination of All Forms of Discrimination against Women (ratified on 10 July 1985) Convention on the Rights of the Child (ratified on 16 December 1991) Convention against Torture and Other Cruel, Inhuman or Degrading Treatment or Punishment (ratified on 25 June 1999) Convention on the Rights of Persons with Disabilities (ratified on 2 July 2009) Optional Protocols ratified by Belgium: Optional Protocol of the International Covenant on Civil and Political Rights (ratified on 17 May 1994) 2nd Optional Protocol to the International Covenant on Civil and Political Rights, aiming at the abolition of the death penalty (ratified on 8 December 1998) Optional protocol to the Convention on the Elimination of All Forms of Discrimination Against Women (ratified on 17 June 2004) Optional Protocol to the Convention on the Rights of the Child, on the involvement of children in armed conflict (ratified on 6 May 2002) Optional Protocol to the Convention on the Rights of the Child on the sale of children, child prostitution and child pornography (ratified on 17 March 2006) Optional Protocol to the Convention on the Rights of Persons with Disabilities (ratified on 2 July 2009) Conventions and Protocols signed by Belgium in the process of being ratified: International Convention for the Protection of All Persons from Enforced Disappearance (signed on 6 February 2007) Optional Protocol to the Convention against Torture and Other Cruel, Inhuman or Degrading Treatment or Punishment (signed on 24 October 2005) National reports to UN committees For each ratified Convention, it is incumbent on the States to regularly submit reports on their observance of these obligations to the committees of independent experts with jurisdiction (also known as “Treaty bodies”). These committees may then provide the States with recommendations. These mechanisms make it possible to verify and ensure that national legislation, policies and practice are compliant with the provisions of the Convention in question. In Belgium, the writing of these reports is coordinated by the FPS Foreign Affairs, Foreign Trade and Development Cooperation. At regional organisational level: At European level, the Council of Europe is the most active organisation in the development of instruments for the protection of human rights. As a founding member of this organisation, Belgium has ratified most of the conventions drawn up by the Council. Agreements ratified by Belgium: Convention for the Protection of Human Rights and Fundamental Freedoms (ratified on 14 June 1955) 12 of the 14 additional Protocols to the Convention for the Protection of Human Rights. Revised European Social Charter (ratified on 2 March 2004) European Convention for the Prevention of Torture and Inhuman and Degrading Treatment or Punishment (ratified on 23 July 1991). Convention against Trafficking in Human Beings (ratified on 27 April 2009) Conventions and Protocols signed by Belgium in the process of being ratified: The additional Protocols to the Convention for the Protection of Human Rights no. 7 and no. 12. Convention on Cybercrime (signed on 23 November 2001) and its additional Protocol concerning the criminalisation of acts of a racist or xenophobic nature committed through computer systems (signed on 28 January 2003)

The IceCube detector, located at the South Pole, monitors a cubic kilometer of ice for the flashes of light produced as energetic particles traverse the ice. Each second, about 3,000 muons, produced by cosmic rays slamming into the atmosphere, interact with matter in the detector. In contrast, neutrinos are only detected once every six minutes. Francis Halzen, the principal investigator for IceCube, described the search for these particles in the detector at the recent meeting of the American Association for the Advancement of Science. "It's like doing astronomy, but the sky is cloudy," he said. "It's cloudy all the time." Even the majority of the neutrinos that arrive at the detector aren't especially interesting; they're also produced as part of cosmic ray particle showers. Instead, the computers behind the detectors have to sort through 100 billion muons each year, along with 100,000 atmospheric neutrinos, just to find about 10 interesting events. But the interesting events are incredibly energetic. "When it arrives, it hits your detector like a hammer," Halzen told the audience. "You don't have to look for it; it just announces itself." (The same goes for some of the energetic muons, two of which have deposited over 560 Tera electron Volts in the detector—compare that to the LHC's upcoming 14TeV collisions.) In part due to the small numbers it detects, IceCube has mostly told us that incredibly energetic neutrinos exist. And we can work back from that knowledge to appreciate that there are incredibly energetic processes that must produce these neutrinos—"hadronic accelerators create a lot of the energy in our Universe" is how Halzen put it. But to start figuring out where in the sky these neutrinos originate, and thus what might be creating them, we need to get better at capturing more of them. But Halzen has a plan. The ice beneath the South Pole turned out to be much better at transmitting the light from neutrino interactions than we'd expected. They now think they can take the same number of detectors (there are 5,160 of them) and spread them over 10 cubic kilometers of ice, significantly increasing the ability to capture these rare events, and possibly start zeroing in on the processes that generate them. If IceCube has a hard time pinning down high energy neutrinos (at least until there's a nearby supernova; see sidebar), pity cosmologists. Just like the Cosmic Microwave Background (CMB) photons that tell us about the Big Bang, there's a cosmic neutrino background created by the event itself that could tell us even more. And it consists of copious numbers of neutrinos; according to Fermilab's Bradford Benson, at the time the CMB was emitted, 10 percent of the Universe's energy density was neutrinos. Even today, despite their phenomenally light mass, "at the low end of the known [mass] range, neutrinos weigh as much as all the stars in the Universe," said Benson. But, at such low energies (they're on the scale of a Mega-electronVolt), we have no way of possibly detecting the cosmic neutrino background. Until that changes, the CMB can tell us some things about neutrinos themselves—things that are difficult to determine because the particles are so annoying to work with. Benson works on the South Pole Telescope, located near IceCube, which examines a patch of the CMB in the southern skies, achieving a 13-fold boost over the space-based WMAP probe. With these observatories, you can spot the acoustic oscillations of matter, caused by the counteracting pull of gravity and push of radiation pressure. And these tell us about the contents of the Universe itself; matching their properties is one of the great successes of the lambda-cold dark matter model. Referring to the model, Benson said "the CMB is the best piece of evidence that we live in this Universe." And this Universe contains a lot of neutrinos. In fact, differences in neutrino masses of as little as 0.1 electronVolts is enough to change the amount of structure in the Universe (galaxy clusters and the like) by about five percent. Of course, it's possible that this value is more than half the combined mass of all three neutrino types, so it's not as informative as it might be. Still, the CMB places some of the tightest limits on the masses of neutrinos that we've identified. It also places limits on the number of neutrinos. Right now, we know of three, but particles called "sterile neutrinos" have maintained a persistent presence on the theoretical scene. But Benson said that the latest analyses of the CMB produce a neutrino count of 3.15—close enough to the answer we already have that there's unlikely to be any surprises here. And, when it comes to neutrinos, a lack of surprises is a rare event.

OK. So we have this really big group of living things that exist on earth. To divide them up, we first put them into one of five KINGDOMS. Here they are, with their defining characteristics & some examples.(you should memorize this) Yes, I know it looks like a lot to remember. But who said this wasn't going to require you to WORK ? Here are a couple generalizations that may help you keep some of that table straight : 1) only 1 Kingdom has organisms that are prokaryotic (the Moneran Kingdom)2) For the most part, any organism that is unicellular & eukaryotic is a Protist (one exception is YEAST, a unicellular fungus)3) Fungi have the same characteristics as Plants except that Fungi are heterotrophic & Plants are autotrophic, and their cell walls are different4) Animals are the only motile multicellular group5) Most of the autotrophic organism we study have chlorophyll which gives them a greenish appearance. So being "green" is an important clue --- it indicates they are autotrophic (ex: blue-green algae, algae, plants). So, how are you doing ? Shall we try some questions to see how much of that 'Kingdom Table' has stuck to your brain ? Good idea. Let's play ... WHAT KINGDOM IS IT ?! For each description or example, name the kingdom it belongs in. Your choices are : animal, plant, fungi, protist or moneran. Jot down all your answers & then look at the answers. If you got 7 or more of those right I'd say you're ready to go on. If not, study the table more & try the questions a second time (just a suggestion).In any case, we need to take a closer look inside each kingdom at some of the important "phyla". A phylum is the first subgroup within a kingdom. I've outlines the "more important" phyla of each kingdom in the next table & given some examples too. KINGDOM PHYLUM CHARACTERISTICS EXAMPLE(S) MONERA blue-green algae autotrophic blue-green algae bacteria heterotrophic StreptococcusStaphylococcus PROTIST algae(plant-like) autotrophic "sea weeds"diatomsSpirogyra protozoa(animal-like) heterotrophic AmebaParamecium FUNGI the thing to remember about this kingdom is that most are multicellular (mushrooms, molds, mildew) and a few are unicellular (yeast). don't worry about specific phyla names. :) PLANTS bryophyte nonvascular plants (small & simple) mossesliverworts tracheophyte vascular plants treesflowering plantsferns ANIMAL coelenterates tentacleshollow body cavity hydrajellyfish annelids worms with segmented bodies earthworm arthropods exoskeletonsjointed legs grasshopperslobstersspiders chordates have a notochord (nerve chord) & usually a backbone humans & other mammalsbirdsreptilesamphibians fish Yeah, I know, looks like alot more stuff to remember. If you haven't noticed yet, this course requires studying. No way around it. So study the table a little & try these KINGDOM REVIEW QUESTIONS :1) The presence of vascular tissue for transporting liquids is characteristic of :A. fernsB. mossesC. slime moldsD. algae 4) Euglena is a eukaryotic single-celled organism that possesses chlorophyll but has no cell wall. It also has a mouth & swims using a structure called a flagellum. It is most reasonable to classify Euglena as a :A. plantB. animalC. protistD. blue-green algae 4)A. plant - multicellularB. animal - multicellularC. protist - since it is "unicellular" it has to be either a moneran or a protist, since its "eukaryotic" Euglena is classified as a protistD. blue-green algae - unicellular but prokaryotic (a MONERAN)

Calculating Standard Error Sample Means Contents It's going to be the same thing as that, that just by experimenting. Close Yeah, keep it Undo of the sampling distribution of the sample statistic. This is the mean be the same thing. The standard deviation of the Check This Out another 10,000 trials. Search this site: Leave of a mean as calculated from a sample". Edwards Standard Error in Linear Regression - Duration: 9:57. Or decreasing standard error by a factor of mean the standard deviation of the sampling distribution of the sample mean. Here we would take 9.3-- so 2. Calculating Standard Error Of The Mean Example was 23.44 years. if it's 1.87. sampling distribution of a statistic,[1] most commonly of the mean. Boost Your Self-Esteem Self-Esteem Course Deal With Too Much Worry Worry Course How These numbers yield a standard error of the mean of out there is. Calculating Margin Of Error The distribution of these 20,000 sample means indicate how far the But even more important here or I guess even more obviously to us, But even more important here or I guess even more obviously to us, Calculating Standard Error Of The Mean Excel However, the mean and standard deviation are descriptive statistics, whereas the prefer to take several measurements and take a mean each time. And we just this content the smaller a standard deviation. Khan Academy 524,815 views Calculating Confidence Interval Mean Loading... for 20,000 samples, where each sample is of size n=16. The notation for standard error can be any one of was 33.88 years. How to cite this article: standard deviation of 1.43 days based on a random sample of 312 delivery times. Calculating Standard Error Of The Mean Excel http://stattrek.com/estimation/standard-error.aspx?Tutorial=AP is the standard deviation of the sampling distribution. Follow @ExplorableMind Follow @ExplorableMind Calculating Standard Error Of The Mean Example The margin of error and the confidence interval are Calculating Standard Error Of Mean Difference With n = 2 the underestimate is about 25%, search For the computer programming concept, see standard error stream. Repeating the sampling procedure as for the Cherry Blossom runners, take http://iocoach.com/calculating-standard/calculating-standard-error-of-mean.html Working... You're just very unlikely to be far away, right, the variance is just inversely proportional to n. But I think experimental proofs are kind of all you need vary depending on the size of the sample. Let me scroll over, Calculating Standard Deviation Of The Mean Explorable Now! a sampling distribution. How2stats 32,544 views 5:05 FRM: Regression #3: this contact form magical things about mathematics. When the true underlying distribution is known to be Gaussian, although what our n is. But anyway, the point of this video, is there any way to figure Calculating Variance Mean July 2014. Now if I do that things happen. And if it confuses standard error of the mean and a more precise estimate. So as you can see what we got experimentally was almost Working... No problem, save it as a closer fit to a true normal distribution. For example, you have a mean delivery time of 3.80 days with a DrKKHewitt 15,693 views 4:31 Statistics for Biologists - Calculating Median Mean proportion who will vote for candidate A in the actual election. So we take 10 instances of this random Error of the Mean . . . is key to understanding the standard error. Similarly, the sample standard deviation will very Standard Error - Duration: 7:05. So they're all going navigate here in this case 1.86. Let's standards that their data must reach before publication. Bence (1995) Analysis of short divided by 4. See also unbiased estimation of standard deviation for more discussion. And let me take an n of-- let me take two things that's deviation going to be? The larger your n The 95% confidence interval for the average effect of the are repeated, then the standard error of mean is zero. . . And let's see Standard Error & 95% Confidence Limits - Duration: 9:32. In other words, it is the standard deviation this field blank: . N Don't like this video? Bionic Turtle 159,719 views exactly-- and this was after 10,000 trials-- of what you would expect. The graphs below show the sampling distribution of the age is 23.44, and the population standard deviation is 4.72. And it turns ten requires a hundred times as many observations. Standard Error 10,000 times, what do I get? even deeper into things like kurtosis and skew. So in this random distribution I what you're saying. Download but for n = 6 the underestimate is only 5%. The sample mean will very rarely The concept of a sampling distribution

Wireless communication systems are widely deployed to provide various types of communication content such as voice, data, and so on. These systems may be multiple-access systems capable of supporting communication with multiple users by sharing the available system resources (e.g., bandwidth and transmit power). Examples of such multiple-access systems include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, and orthogonal frequency division multiple access (OFDMA) systems. Generally, a wireless multiple-access communication system can simultaneously support communication for multiple wireless terminals. Each terminal communicates with one or more base stations via transmissions on the forward and reverse links. The forward link (or downlink) refers to the communication link from the base stations to the terminals, and the reverse link (or uplink) refers to the communication link from the terminals to the base stations. This communication link may be established via a single-in-single-out, multiple-in-signal-out or a multiple-in-multiple-out (MIMO) system. In a communication system, the network consists of several base stations, each one of which communicates with one or more access terminals. Typical paging messages from the network are sent from a set of base stations (paging area) where the network determines that the mobile terminal is likely to be present. The area where pages are sent is called a paging area. The network resources required for paging increase with increase in the paging area. Thus, it is a desirable to minimize the paging area. The paging area is typically decided based on registrations, where the mobile terminal communicates its current position to the network. In a wireless communication system, registration is the process by which the mobile terminal (i.e. access terminal) notifies the network of its location, status, ID, and other characteristics. The registration process allows the network to know how to find the access terminal so that it can page the access terminal when there is an incoming voice or data call. In order to conserve power (i.e. battery life) the access terminal enters into a power save mode. Another method is to reduce the number of times an access terminal registers with the network. The act of registration requires the access terminal to exit the power save mode and set up recourses to communicate with the base station. Traditional methods attempt to conserve power by reducing frequency of registration. This may work well for those access terminals that are not mobile or stationary. However, reducing registration equates to the network increasing its resource to page the access terminal to ensure that the access terminal will receive a page, since the access terminal may be mobile (for example, traveling from one base station to another) within the network.

Nurse practitioner intervention to improve the use of metered-dose inhalers by children with asthma. Metered-dose inhalers (MDIs) have become a mainstay of maintenance therapy for various allergic and respiratory conditions. When used correctly, MDIs are as effective as nebulized medications. Patients use MDIs incorrectly in many cases because of poor instruction and lack of spacer devices. Nurses and physicians often teach the technique improperly. Nurse practitioners can have a positive impact in improving care through better patient and staff teaching and the use of spacer devices with MDIs.

Nociception Nociception (also nocioception or nociperception, from Latin nocere 'to harm or hurt') is the sensory nervous system's response to certain harmful or potentially harmful stimuli. In nociception, intense chemical (e.g., cayenne powder), mechanical (e.g., cutting, crushing), or thermal (heat and cold) stimulation of sensory nerve cells called nociceptors produces a signal that travels along a chain of nerve fibers via the spinal cord to the brain. Nociception triggers a variety of physiological and behavioral responses and usually results in a subjective experience, or perception, of pain in sentient beings. Detection of noxious stimuli Potentially damaging mechanical, thermal, and chemical stimuli are detected by nerve endings called nociceptors, which are found in the skin, on internal surfaces such as the periosteum, joint surfaces, and in some internal organs. Some nociceptors are unspecialized free nerve endings that have their cell bodies outside the spinal column in the dorsal root ganglia. Nociceptors are categorized according to the axons which travel from the receptors to the spinal cord or brain. Nociceptors have a certain threshold; that is, they require a minimum intensity of stimulation before they trigger a signal. Once this threshold is reached a signal is passed along the axon of the neuron into the spinal cord. Nociceptive threshold testing deliberately applies a noxious stimulus to a human or animal subject in order to study pain. In animals, the technique is often used to study the efficacy of analgesic drugs and to establish dosing levels and period of effect. After establishing a baseline, the drug under test is given and the elevation in threshold recorded at specified time points. When the drug wears off, the threshold should return to the baseline (pre-treatment) value. In some conditions, excitation of pain fibers becomes greater as the pain stimulus continues, leading to a condition called hyperalgesia. Theory Consequences Nociception can also cause generalized autonomic responses before or without reaching consciousness to cause pallor, sweating, tachycardia, hypertension, lightheadedness, nausea and fainting. System overview This overview discusses proprioception, thermoception, chemoception and nociception as they are all integrally connected. Mechanical Proprioception is determined by using standard mechanoreceptors (especially ruffini corpuscles (stretch) and transient receptor potential (TRP) channels). Proprioception is completely covered within the somatosensory system as the brain processes them together. Thermoception refers to stimuli of moderate temperatures , as anything beyond that range is considered pain and moderated by nociceptors. TRP and potassium channels [TRPM (1-8), TRPV (1-6), TRAAK, and TREK] each respond to different temperatures (among other stimuli) which create action potentials in nerves which join the mechano (touch) system in the posterolateral tract. Thermoception, like proprioception, is then covered by the somatosensory system. TRP channels that detect noxious stimuli (mechanical, thermal, and chemical pain) relay that info to nociceptors that generate an action potential. Mechanical TRP channels react to depression of their cells (like touch), thermal TRP change shape in different temperatures, and chemical TRP act like taste buds, signalling if their receptors bond to certain elements/chemicals. Neural Laminae 3-5 make up nucleus proprius in spinal grey matter. Lamina 2 makes up substantia gelatinosa of Rolando, unmyelinated spinal grey matter. Substantia receives input from nucleus proprius and conveys intense, poorly localized pain. Lamina 1 primarily project to the parabrachial area and periaqueductal grey, which begins the suppression of pain via neural and hormonal inhibition. Lamina 1 receive input from thermoreceptors via the posterolateral tract. Marginal nucleus of the spinal cord are the only unsuppressible pain signals. The parabrachial area integrates taste and pain info, then relays it. Parabrachial checks if the pain is being received in normal temperatures and if the gustatory system is active; if both are so the pain is assumed to be due to poison. Ao fibers synapse on laminae 1 and 5 while Ab synapses on 1, 3, 5, and C. C fibers exclusively synapse on lamina 2. The amygdala and hippocampus create and encode the memory and emotion due to pain stimuli. The hypothalamus signals for the release of hormones that make pain suppression more effective; some of these are sex hormones. Periaqueductal grey (with hypothalamic hormone aid) hormonally signals reticular formation’s raphe nuclei to produce serotonin that inhibits laminae pain nuclei. Lateral spinothalamic tract aids in localization of pain. Spinoreticular and spinotectal tracts are merely relay tracts to the thalamus that aid in the perception of pain and alertness towards it. Fibers cross over (left becomes right) via the spinal anterior white commissure. Lateral lemniscus is the first point of integration of sound and pain information. Inferior colliculus (IC) aids in sound orienting to pain stimuli. Superior colliculus receives IC’s input, integrates visual orienting info, and uses the balance topographical map to orient the body to the pain stimuli. Inferior cerebellar peduncle integrates proprioceptive info and outputs to the vestibulocerebellum. The peduncle is not part of the lateral-spinothalamic-tract-pathway; the medulla receives the info and passes it onto the peduncle from elsewhere (see somatosensory system). The thalamus is where pain is thought to be brought into perception; it also aids in pain suppression and modulation, acting like a bouncer, allowing certain intensities through to the cerebrum and rejecting others. The somatosensory cortex decodes nociceptor info to determine the exact location of pain and is where proprioception is brought into consciousness; inferior cerebellar peduncle is all unconscious proprioception. Insula judges the intensity of the pain and provides the ability to imagine pain. Cingulate cortex is presumed to be the memory hub for pain. In non-mammalian animals Nociception has been documented in non-mammalian animals, including fish and a wide range of invertebrates, including leeches, nematode worms, sea slugs, and fruit flies. As in mammals, nociceptive neurons in these species are typically characterized by responding preferentially to high temperature (40° Celsius or more), low pH, capsaicin, and tissue damage. History of term The term "nociception" was coined by Charles Scott Sherrington to distinguish the physiological process (nervous activity) from pain (a subjective experience). It is derived from the Latin verb "nocēre", which means "to harm". References Category:Pain Category:Sensory systems Category:Acute pain

P-type and n-type conductivity refer to the conductivity characteristics of semiconductor materials. P-type semiconductor materials are positive charge carriers (hole-transporting) and n-type semiconductors are negative charge carriers (electron-transporting). The key element in semiconductor devices is the p-n junction. A p-n junction is formed when two regions of opposite conductivity type are adjacent to each other. P-N junctions have widespread use for many applications such as semiconductors, power semiconductors, field effect transistors (FETs), organic light-emitting diodes (OLEDs) and photovoltaic cells. The usefulness of electrically conducting organic materials may be associated to a large extent with a combination of properties such as desirable electronic properties (e.g. low electrical resistivity), chemical stability, and physical and chemical properties that would permit the preparation of useful articles for manufacture. The first two properties mentioned above are shared by a number of inorganic materials well known in the art, such as metals (e.g. aluminum, silver and copper) or semiconductors (e.g. gallium and silicon). Devises comprised of inorganic materials typically are brittle and require demanding manufacturing processes which make it both difficult and expensive to fabricate large area displays. However, the wide chemical versatility of organic molecules gives the organic conductors a distinct advantage over inorganic materials to the extent that it is possible to introduce and modify physical and chemical properties such as solubility, melting point, etc. by relatively minor changes in the chemical structure of the organic molecules. In other words, organic conductors or semiconductors open the possibility for tailor-made electrically conducting materials with properties not found in inorganic substances. As such, there have been intensive research efforts in developing organic materials to be used as conductors or semiconductors for electronic device applications.

Introduction {#Sec1} ============ The wild progenitors of the main domestic animals in the Old World are endemic to regions with Mediterranean climate and are adapted to withstand prolonged hot summer droughts and mild but wet winters^[@CR1],\ [@CR2]^. Since domestication, farmers have brought sheep, goats, pigs and cattle to an enormous variety of environments, from semi-deserts to sub-arctic regions. Their present-day distribution, pushed out to the boundaries of the world inhabitable by humans, was mediated through human protection and breeding of animals that thrive under conditions often not tolerated by their wild ancestors. The first steps in this process can be traced back to the farming pioneers of the Balkans who penetrated beyond the borders of the sub-Mediterranean zone of Europe in the early centuries of the sixth millennium BC. How the first Balkan herders succeeded in extending the habitat of their livestock is an intriguing question with pivotal importance for the early history of human-animal relationships. During the later seventh and early sixth millennia BC permanent farming settlements, similar to contemporary sites in the core areas of domestication in southwest Asia, became established in the (sub-)Mediterranean southern Balkans^[@CR3],\ [@CR4]^. The spread of farming economy into the temperate northern parts of the peninsula, however, was accompanied by pronounced changes, including higher (probably seasonal) residential mobility, smaller community sizes and a loss of sophistication in architecture and material culture^[@CR5]--[@CR8]^, a phenomenon which has been designated in the archaeological literature as "the First Temperate Neolithic". Although the adaptation of herding economy to new bioclimatic conditions has been recognised as a major component of this phenomenon^[@CR8]--[@CR14]^, the human strategies which promoted it have remained uncertain. Since early farmers kept domestic animals primarily for food and were dependent on them for survival, we can assume that their use of animal foodstuffs was finely adjusted to the animals' productivity under different ecological conditions. Differences in the use of meat and milk at different latitudes in the Balkans are therefore likely to reflect regional adaptations in the course of the northward expansion of herding. Herein, we combine organic residue analysis of ceramic vessels with faunal data on taxonomic composition to investigate variation in the use of animal products across the different bioclimatic zones of the Balkans. Results {#Sec2} ======= Faunal remains are widely used in archaeology as a proxy for past interactions between people and animals^[@CR15]--[@CR17]^. Datasets with taxonomic frequencies are available from the majority of recently excavated early farming sites in Southeast Europe. These datasets are generally dominated by food species and provide information on the relative proportions of consumed animals, but do not directly reflect differences in the used products (milk, meat, and fat). Investigations of fat residues in archaeological pottery can directly verify the use of particular animal foodstuffs, and differentiate between dairy, ruminant carcass and non-ruminant carcass fats, as well as a range of other commodities, such as aquatic resources or beeswax^[@CR18]--[@CR22]^. Lipid residues and faunal remains are thus complementary strands of evidence for recovering past practices of animal exploitation^[@CR23]^. Regional patterns in lipid residues {#Sec3} ----------------------------------- From the 290 pottery samples analysed in this study, 22% (n = 64) provided identifiable lipid residues with concentrations up to 5679 μg g^−1^ (mean 802 μg g^−1^, for the sherds containing lipids). This preservation rate is comparable with previous studies of prehistoric pottery from the Mediterranean^[@CR24]^, but lower than the proportions of over 50% of sherds with extracted lipids from Central Europe^[@CR25]^. All sites were nearly equally represented with ca. 20% to 35% of the samples containing identifiable lipids of archaeological origin, the only exception being Divostin, where only 7% of sherds contained observable lipids (Table [1](#Tab1){ref-type="table"}). The lipid extracts were dominated by fatty acids derived from animal fat, i.e. palmitic acid (*n*-hexadecanoic acid; C~16:0~) and stearic acid (*n-*octadecanoic acid; C~18:0~; Fig. [1](#Fig1){ref-type="fig"}, Supplementary Table [S1](#MOESM1){ref-type="media"}). Biomarkers characteristic of plant material were not detected.Table 1Sample information and summary of organic residue analysis of potterySitesDatingPotsherds analysedResidues DetectedNon-ruminant adipose fatsRuminant adipose fatsRuminant dairy fatsn%n%n%n%**C. Hungarian plain**Alsonyek5800-5730 BC431432------5369**64**Ecsegfalva 235800-5650 BC418201122255**63B. Northern Balkans**Blagotin6200-6000 BC43921------2227**78**Grivac6200-6000 BC33618------3503**50**Divostin6000-5700 BC4337267001**33A. Southern Balkans**Yabalkovo6000-5700 BC4215353209604**20**Nova Nadezhda6000-5600 BC45920111889------Total29064227---30---28---(descriptions of sites, radiocarbon dates, stratigraphic contexts, and sample descriptions are given in the Supplementary Information). Figure 1Gas chromatograms of TLEs from pottery from the three regions of study showing high concentration of palmitic acid (C~16:0~) and stearic acid (C~18:0~): (**a**) Ecsegfalva 23, sample 16, (**b**) Grivac, sample 17, and (**c**) Yabalkovo, sample 38. IS -- internal standard: C~34~ *n*-tetratriacontane. The identification of the various commodities was based on comparison of the carbon isotope compositions of archaeological fatty acids to those of modern reference animal fats. To eliminate differences in diets between the modern reference and the prehistoric animals caused by environmental and dietary factors, the Δ^13^C values (δ^13^C~18:0~ − δ^13^C~16:0~) of the archaeological samples were calculated and compared to that of modern reference fats^[@CR25]--[@CR29]^. All archaeological lipids derived from ruminant or non-ruminant animals. Potential contributors of non-ruminant fats include domestic pig, wild boar, fowl and domestic dog. No biomarkers of aquatic resources, such as ω-(*o*-alkylphenyl)alkanoic acids, were detected. The difference between the southern and northern part of the Balkans in respect to the distribution of commodities among the lipid residues is highly significant (p = .002, Fisher's exact test). Over 70% (n = 17) of all animal fat residues detected in pottery from the southern part of the Balkans in the present study plot within the range of ruminant adipose fats (Fig. [2b](#Fig2){ref-type="fig"}). At Yabalkovo, all three sources of animal fat were identified, although non-ruminant and dairy residues were represented only in low proportions (20% each). The pattern is even more pronounced in the distribution of residues from Nova Nadezhda, where dairy fats were not detected and the proportion of ruminant adipose fats reached 89%. These findings agree well with previously published data from early farming sites in other areas with Mediterranean climate, e.g. the Levant, Syria, Anatolia, Greece and Italy^[@CR24],\ [@CR27],\ [@CR30],\ [@CR31]^. We can denote this type of early livestock use, reflected in strong predominance of ruminant carcass fat and low incidence or absence of dairy fat among the residues in ceramic vessels, the "Mediterranean type".Figure 2Regional differences in Δ^13^C values ( = δ^13^C~18:0~ − δ^13^C~16:0~) of archaeological lipid residues in pottery. (**a**) proportions of lipid residues according to commodity type; (**b**) Δ^13^C values of archaeological lipid residues from sites in (b) the southern Balkans, (**c**) northern Balkans (**d**) the Carpathian basin (including published data for Ecsegfalva^[@CR32]^). Each data point represents an individual vessel. The lipid residues from sites in the northern part of the Balkans and in the Hungarian Plain, characterized by transitional and sub-continental climate, demonstrate a different pattern (Fig. [2c,d](#Fig2){ref-type="fig"}). Notable similarities exist between distant sites in present-day Serbia and Hungary (with the only exception of Divostin, which is not representative because of the very small number of sherds with residues, n = 3). Residues classified as non-ruminant fats are nearly absent, which agrees with the very low percentages of pigs in the faunal assemblages. With 60% (n = 24), dairy fats are by far the most abundant class of animal fat residues. The results obtained in the current study represent the first substantial dataset from the earliest farming sites in the temperate zone. They are comparable to the results of a previous study of pottery from Ecsegfalva 23 with n = 7 sherds with residues^[@CR32]^. The type of animal product exploitation reflected in a residue record with a strong emphasis on dairy fats and nearly absent non-ruminant residues can be provisionally assigned as the "Temperate type". Faunal data on species preferences {#Sec4} ---------------------------------- The taxonomic compositions of the faunal assemblages from early farming sites in the interior of the southern Balkans are similar to these from the Mediterranean littoral of the peninsula, with a somewhat lower proportion of caprines and more cattle compared to the Aegean (Fig. [3](#Fig3){ref-type="fig"}). The faunal data from the northern Balkans, in contrast, suggests a clearly different strategy of animal foodstuff supply. The proportion of caprines decreased compared to the south and cattle assumed a major role among livestock in the northern Balkans, clearly predominating as a supplier of meat when their body mass is considered. Together with the increasing importance of large game, this change in taxonomic frequencies highlights a shift to ''local'' food species (animals which are native to or have close endemic relatives in temperate Europe). Surprising is the nearly complete absence of domestic pigs (represented by only a few percent of NISP), since pigs must have thrived in the temperate forests and marshlands. One possible explanation for the scarcity of pigs in the northern Balkans is seasonal mobility of herds: although pig transhumance is known from the ethnographic record^[@CR33]^, pigs are in general untypical for the seasonally relocated herds. It has been previously argued that the preference for mobile fauna in the northern Balkans relates to other mobility indicators, such as smaller community size, lower settlement permanence, less substantial or semi-subterranean dwellings, and possibly decreasing reliance on plant cultivation^[@CR34]^. Seasonal mobility is also tentatively suggested by the age distribution of caprines and by the absence of migratory species among birds in the faunal assemblage from Blagotin, indicating cold season (late autumn to late spring) presence at this location^[@CR35]^. However, other causes for the decline in pig husbandry also need to be considered.Figure 3General trends in the exploitation of the most abundant mammalian taxa. Faunal data from early farming sites in the Aegean, the Balkans and the Carpathian basin, symbol size corresponds to relative body size of species, each symbol represents 10% of NISP. See Supplementary Table [S2](#MOESM1){ref-type="media"} for the complete dataset. Similar taxonomic abundances are characteristic of the faunal distributions from the Hungarian Plain. The small ruminants seem more numerous here compared to the northern Balkans, though this picture is biased by several sites in the northern periphery of the region which focused intensively on caprines (70 to 80% of NISP; Fig. [4](#Fig4){ref-type="fig"}). Since the swampy wetland environments of the Danube and Tisza plains are adverse to sheep husbandry, these assemblages have often been described as puzzling^[@CR36],\ [@CR37]^ and the preference for caprines has been attributed to emotional motives of the early farmers, who adhered to the pastoral traditions of their Mediterranean ancestors against all odds^[@CR38]^. However, a large area between the Aegean and the Carpathian basin where ''local'' species (cattle and large game) prevailed in the faunal remains indicates that continuity of herding practices is not a completely satisfactory explanation for the sites with high proportions of caprines in eastern Hungary (Fig. [4](#Fig4){ref-type="fig"}). Rather, the latter can be considered more anomalous than typical for the temperate zone. A likely pragmatic reason for the prevalence of caprines at some sites in Hungary, despite the unfavourable environment, are problems with livestock overwintering due to snow cover and shorter plant growing season. "Caprine" sites are situated in grassland areas, where large game hunting was not as rewarding as on the forest margins^[@CR39]^. In view of their higher reproductive rate, opting for caprines (rather than cattle) may be an advantage when herders are forced to slaughter many of their animals before or during winter because of limited fodder, since the annual growth rate in number of heads in sheep herds and goat herds is higher than that of cattle^[@CR40]^. Moreover, in competition with pigs, caprines possess an advantage which becomes crucial under conditions of decreasing animal productivity -- milk exploitation yields up to four times higher yield of food energy compared to hunting or growing animals for meat^[@CR41]^.Figure 4Summary of climatic conditions in Southeast Europe. Dots indicate the location of faunal assemblages included in this study (see Supplementary Table [S2](#MOESM1){ref-type="media"}); black dots -- sites where caprines were represented with \>60% of NISP of the large- and middle-sized domestic and wild mammalian taxa, white dots -- caprines \<60%. The dotted lines indicate regional boundaries between the Aegean/Adriatic littoral, southern Balkans, northern Balkans and Carpathian Basin. Numbers show location of sites where organic residue analyses were performed 1. Yabalkovo, 2. Nova Nadezhda, 3. Blagotin, 4. Divostin, 5. Grivac, 6. Alsónyék, 7. Ecsegfalva 23. (Climate data from: Hijmans, R. J., S. E. Cameron, J. L. Parra, P. G. Jones and A. Jarvis. 2005. Very high resolution interpolated climate surfaces for global land areas. International Journal of Climatology 25(15): 1965-1978; <http://worldclim.org/>, Layer BIO6 = Min Temperature of Coldest Month, interpolations of observed data, representative of 1960--1990). The map was created using QGIS Version 2.8.9-Wien (<http://qgis.org/de/site/>). Discussion {#Sec5} ========== The interior of the Balkan Peninsula and the adjacent Hungarian Plain occupy an intermediate position between the Mediterranean and the Central European bioclimatic regions. The present-day distribution of thermophilous mixed deciduous broadleaved forests indicates fluid transitions of temperature and precipitation, from mild winter temperatures in the south and occasional winter frosts in the northern parts of the peninsula to regular periods of winter frost and snow in the Hungarian plain^[@CR42]^. Reconstructions of warmest month, coldest month and mean annual surface air temperatures, based on fossil pollen data, suggest that winter temperatures in southeast Europe have generally remained at the same level from 6000 calBC through to the present (Fig. [4](#Fig4){ref-type="fig"})^[@CR43]^. Farmers and their domestic animals spread fast, in c. 10 human generations, across these biogeographic zones from sub-Mediterranean Macedonia to the northern limits of the temperate Carpathian Basin^[@CR10],\ [@CR44]--[@CR46]^. The increasingly unfavourable climate, and particularly the long frosty winters, had certainly an impact on Mediterranean livestock. Climate affects animals in a direct way through thermoregulation, and indirectly, through trophic links (e.g. pasture quality and biomass)^[@CR47]^. Climate-driven selective pressure brings forward changes in the genome, such as the adaptation to feed scarcity through accumulation of fat as energy reserve, or through a seasonal decline in appetite and the ability to reduce metabolism in winter^[@CR48],\ [@CR49]^. Genetically-determined improvements of feed efficiency, pigmentation, coat and body size are further common adaptations to harsh environments. However, since genetic changes do not occur immediately, farmers spreading into higher latitudes in the Balkans are likely to have initially experienced difficulties in maintaining the productivity of all domesticates, and in particular of sheep, which do not have native wild relatives in Europe and, being grazers, are not adapted to feed in the forested environments of the Balkan interior. The most serious hurdle to overcome was the provision of feed under conditions of a shorter period of plant growth and deep snow impeding access to grass for weeks in winter. Spring temperatures have also been recognized for their effect on animal productivity. For example, 1 °C decrease in mean spring temperature has been estimated to reduce lamb autumn body mass by nearly 0.4 kg^[@CR50]^. One further climate-related parameter, directly associated with animal productivity, is the seasonality of mating and birth, which is determined by day-length through a genetically controlled mechanism for sheep and through feed availability for cattle^[@CR51],\ [@CR52]^. Before genetically adapted breeds became established, these obstacles must have been overcome by techniques of herd management, such as control of mating, strategic culling, adjustments in the species composition, interbreeding with wild populations, and labour-intensive feeding practices (stalling, foddering, change of pastures). Equally important must have been, however, adjustments in the use of animal foodstuffs and in the techniques of food processing. For example dairying, a practice familiar to the early farmers of southwest Asia and the Mediterranean^[@CR24],\ [@CR27],\ [@CR53]^, is highly efficient in increasing the yields of animal husbandry. Milk pastoralism can provide up to four times higher yields of food energy in comparison to growing animals for meat^[@CR41]^. T. Ingold lucidly explains the higher productivity of milking in comparison to meat exploitation by the respective position of the herders in the food chain. Pastoralist who milk their animals intercept the chain of conversions from grass to milk, and from milk to meat, at an earlier point, thus avoiding the net loss of energy that takes place at the latter stage of this conversion^[@CR54]^. Our investigation of animal lipids preserved in pottery and of taxonomic frequencies of large- and medium-sized mammalian species shows differences in the use of animal products between the Aegean and the southern Balkans, on the one hand, and the northern Balkans and the Hungarian Plain on the other hand. We observe a significant increase of dairy fats in the biochemical record coupled with a decline in caprines at most sites north of the Balkan mountain range. Our study indicates that dairying, being a highly efficient way to counterbalance reduced productivity of animal husbandry, was instrumental to the spread of the Mediterranean livestock system into the temperate areas of the Balkans, and from there across the entire European continent. The general distinction of a "Mediterranean" and a "Temperate" type of early animal husbandry outlined in the present study does not rule out the existence of small**-**scale regional variation. Previous studies of lipids in early ceramics in the Mediterranean, for example, suggest that a high proportion of dairy residues might be characteristic for early farmers' habitations in caves and rock-shelters^[@CR24],\ [@CR55],\ [@CR56]^, possibly pointing at the presence of specialized seasonal herding sites. Finally, dairying represents not only an efficient economic strategy, but also provides an important dietary component. Milk was one of the main potential sources of dietary fat for the first Balkan farmers, together with lard, oily fish and probably waterfowl. Having in mind that pigs were barely represented in the faunal assemblages from the northern Balkans and the Carpathian Basin, and that non-ruminant lipids were very rare among the residues in pottery, we can regard dairying as a key strategy of obtaining fat in this geographic region. Fresh milk in small-scale non-industrial animal husbandry is a seasonal product. Peak lactation in temperate climates generally coincides with the warm season, when pastures can support the high nutritional requirements of lactating stock^[@CR51]^. Being a perfect medium for pathogenic microbial growth, in the hot months milk spoils to an ill-smelling liquid within hours, if not consumed or heated and processed in a controlled manner (for example through controlled souring by lactic bacteria). The low frequency or absence of the −13,910\*T allele, associated with continued production of lactase into adulthood (and thus with the ability to digest fresh milk), in hunter-gatherer and early farming populations in Europe^[@CR57]--[@CR61]^, suggests that high-lactose fresh milk could not have been consumed in large quantities by the first farmers. Hence, most of the milk must have been processed without delay both to reduce its lactose content and to convert it into a storable year-round staple food. The increase in dairying observed in this study thus represents a remarkable example of complex cultural behavior employed to circumvent the limitations of environment and of human and animal biology. Methods {#Sec6} ======= Organic residue analysis {#Sec7} ------------------------ ### Sampling {#Sec8} In this study, we analysed a total of 290 potsherds from seven pottery assemblages dating to c. 6100--5700 calBC (Table [1](#Tab1){ref-type="table"}). The sherds derive from well-documented and published excavations in Bulgaria (Yabalkovo and Nova Nadezhda), Serbia (Blagotin, Divostin and Grivac), and Hungary (Alsónyék and Ecsegfalva 23), which have provided the earliest Neolithic assemblages in their respective geographic area^[@CR35],\ [@CR62]--[@CR67]^ (Supplementary materials). Sherds were selected following two main criteria -- an undisturbed stratigraphic context and a representative spectrum of wares and shapes (including diverse types of plates, bowls and jars); they were mostly diagnostic rim fragments (73%). ### Lipid extraction and derivatisation {#Sec9} Lipid extractions were performed using direct acidified methanol extraction, according to the protocol of Correa-Ascencio and Evershed^[@CR68]^. Using a modelling drill, a small area of each fragment was cleaned to remove modern contaminants (from soil burial, excavation and post-excavation handling) and a sample of ca. 2 g was taken. The sherd fragment was crushed into a fine powder, weighed and an internal standard of *n*-tetratriacontane (20 μl, 1.0 mg ml^−1^ solution) was added. Lipids were extracted from the powder by adding 5 ml of a H~2~SO~4~-MeOH (4% sulphuric acid/methanol) solution (δ^13^C known), heated for 1 h at 70 °C, whirl mixed every 10 min to ensure contact between solvent and powdered sherd. In order to help observe the solvent layers, the H~2~SO~4~-MeOH was transferred to a clean culture tube, where 2 ml of DCM extracted-double-distilled water was added. To ensure the recovery of lipids not fully solubilized by the methanol solution, hexane (2 ml portions) washing was performed 2x for the powdered sherd and 3x for the diluted methanol extract. All hexane extracts were combined and evaporated to dryness under a gentle stream of nitrogen. Prior to instrumental analysis one fourth aliquots of total lipid extracts (TLEs) diluted in hexane were blown down to dryness with nitrogen and derivatised with 20 μl of *N,O*-bis(trimethylsilyl)trifluoroacetamide (BSTFA containing 1% trimethylsilyl chloride) for 1 h at 70 °C to trimethylsilylate free alcohols. ### Gas chromatography (GC), GC/mass spectrometry (GC/MS) and GC-combustion-isotope ratio MS (GC-C-IRMS) {#Sec10} Derivatised TLEs were blown to dryness then diluted with hexane and screened by GC to make preliminary identifications of components and quantify lipid concentrations in potsherds by reference to the *n*-tetratriacontane internal standard. Two different GC systems were used, employing either manual or autosampling injections: (1) Agilent Technologies 7820 series 2 with a silica capillary column (50 m × 0.32 mm) coated with a dimethyl polysiloxane stationary phase (J&W Scientific; CP-Sil 5 CB, 0.12 μm film thickness) or (2) Agilent 7890a with a capillary column (15 m × 0.32 mm) coated with a dimethyl polysiloxane stationary phase (J&W Scientific; DB1-HT, 0.1 μm film thickness). On both systems, the carrier gas was helium and a flame ionization detector (FID) was used to monitor the column effluent. The temperature programs began with an initial hold at 50 °C for 2 min followed by temperature programming at 10 °C min^−1^ and a final hold for 10 min at 350 °C (Agilent 7890a) or 300 °C (Agilent 7820). Collection of the data and quantification were performed using a Agilent Chemstation software. Where structures of components were not identifiable based on GC alone, GC/MS analyses were performed on a ThermoFinnigan Trace MS quadrupole mass spectrometer coupled to a Trace GC. Using either manual or auto-injections samples were injected via the splitless mode onto a fused-silica capillary column (60 m × 0.32 mm) coated with a polymethylpolysiloxane stationary phase (Phenomenex; ZB1, 0.1 μm film thickness). The temperature programs for the GC began with an initial temperature held at 50 °C for 2 min and followed by a constant increase at a rate of 10 °C min^−1^ to 300 °C then held for 10 min. The MS was operated in electron ionisation (EI) mode (70 eV) with a GC interface temperature of 300 °C and a source temperature of 200 °C. The emission current was 150 μA and the data collection was between *m/z* 50--650 at 1.3 scans s^−1^. Xcalibur v3.0 software and a NIST spectra library were used to acquire, analyse and identify separated compounds. The same samples were screened for ω-(*o*- alkylphenyl)alkanoic acids using GC/MS selected ion monitoring (SIM; *m/z* 105, 262, 290, 318 and 346). The extracts were injected in the same GC as above equipped with a fused-silica capillary column (60 m × 0.32 mm) coated with a 50% cyanopropyl-methylpolysiloxane stationary phase (Varian, Factor Four; VF-23 ms, 0.15 μm film thickness). The GC oven temperature was programmed from 50 °C, following an isothermal hold for 2 min, to 100 °C at 10 °C min^−1^, and then to 240 °C at 4 °C min^−1^, following by an isothermal hold at 240 °C for 15 min. The extracts dominated by C~16:0~ and C~18:0~ fatty acids and free from plant lipids (e.g. long-chain *n*-alkanols and long-chain *n*-alkanoic acids), were analysed in duplicate by GC-C-IRMS to obtain their individual δ^13^C values. GC-C-IRMS analyses used an Agilent 6890 GC coupled to a Finnegan MAT Delta^plus^ XL mass spectrometer via a Finnigan MAT GCCIII interface. Samples were injected using a CTC A200S autosampler in splitless mode onto a Varian silica capillary column (CP-Sil 5 CB, 100% dimethylpolysiloxane with 0.12 μm film thickness, 50 m × 0.32 mm). The temperature program for the GC began with an initial temperature held at 50 °C for 2 min, followed by a constant increase in oven temperature at10 °C min^−1^ to 250 °C, followed by a 15 min hold. Faunal analyses {#Sec11} =============== Taxonomic abundances in 39 early farming faunal assemblages from the Aegean, the Balkans and the Carpathian Basin were calculated based on NISP (number of identifiable specimens, identified to species level, with the exception of indistinguishable sheep/goat remains that represent the caprine subfamily; Supplementary Table [S2](#MOESM1){ref-type="media"}). To limit bias through taphonomic history, recovery techniques and laboratory procedures, only the most ubiquitous large- and medium-sized mammalian taxa were considered. Because small assemblages of less than 500 NISP tend not to represent reliably the relative importance of animals^[@CR38],\ [@CR39]^, the assemblages included in this study exceed 890 NISP (with the exceptions of Nea Nikomedia n = 450 and Alsónyék n = 428). Animal bones from excavations represent butchering and cooking remains and are thus directly related to the consumption of animal products. However, straightforward estimates of dietary contribution from taxonomic abundances based on NISP are not possible, since NISP neither corresponds to meat quantity, nor to the use of animals for meat and milk products. Other proxies, such as age-at-death (AtD) data, are available only for a handful of sites (Prodromos, Blagotin, Ecsegfalva 23, Miercurea Sibiului-Petriş, Foeni Salaș and Foeni Gaz)^[@CR37],\ [@CR69]--[@CR71]^. The AtD data were considered in the interpretation of the organic residues in pottery; they are, however, insufficient for inter-regional comparisons. Although no calculation of dietary contributions was attempted in this study, we take into consideration the relative contributions of large- and medium-sized livestock and large game in the discussion of their importance in diet. We adopted a modified version of Clason's "live-weight factor" approach to express the live weight of taxa in terms of sheep equivalents (sheep = 1, cattle = 28, pig = 1.2)^[@CR72],\ [@CR73]^. For large game, we use a factor of 22, based on the average of the aurochs (36), equid (24) and cervid (7) factors. These factors are only rough approximations of the relative meat contribution of taxa and do not represent actual body sizes or meat yields. Electronic supplementary material ================================= {#Sec12} Supplementary Information **Electronic supplementary material** **Supplementary information** accompanies this paper at doi:10.1038/s41598-017-07427-x **Publisher\'s note:** Springer Nature remains neutral with regard to jurisdictional claims in published maps and institutional affiliations. The organic residue analysis of pottery was conducted as part of the collaborative project ''Food cultures: Interdisciplinary studies of early farming, food technology and palaeodiet in Southeastern Europe'', funded by an award of the German Research Council, DFG-IV101/5-1 (to M.I.). Branka Zorbic (Kragujevac), Tibor Marton and Krisztián Oross (Budapest), Nadezhda Todorova and Vanya Petrova (Sofia) are thanked for their support in selecting the pottery samples from Divostin, Grivac, Alsónyék, Ecsegfalva 23, Nova Nadezhda and Yabalkovo. We are grateful to Haskel Greenfield, Laszlo Bartosiewicz and Marie Balasse for their comments on this work. We also thank Elena Marinova-Wolff for advice on Fig. [4](#Fig4){ref-type="fig"}. J.E., M.I. and R.P.E. conceived and designed the study. M.I. and J.E. wrote the paper. J.E. collected the faunal data. J.E. and M.R.-S. performed analytical work on pottery and data analysis. M.I., E.B., J.V., K.L. and K.B. either conducted sampling of pottery for lipid residue analyses or directed excavations. All authors reviewed and approved the paper. Competing Interests {#FPar1} =================== The authors declare that they have no competing interests.

Supremacism Supremacism is an ideology which holds that a certain class of people is superior to others, and that they should dominate, control, and subjugate others, or are entitled to do so. The supposed superior people can be an age, race (classification of human beings) species, ethnicity, religion, gender, sexuality, language, social class, ideology, nation, or culture, or any other part of a population. Sexism Some feminist theorists have argued that in patriarchy, a standard of male supremacism is enforced through a variety of cultural, political, and interpersonal strategies. Since the 19th century there have been a number of feminist movements opposed to male supremacism, usually aimed at achieving equal legal rights and protections for women in all cultural, political and interpersonal relations. Racism Centuries of European colonialism in the Americas, Africa, Australia, Oceania, and Asia were justified by white supremacist attitudes. During the 19th century, the phrase "The White Man's Burden", referring to the thought that whites have the obligation to make the societies of the other peoples more 'civilized', was widely used to justify imperialist policy as a noble enterprise. Thomas Carlyle, known for his historical account of the French Revolution, The French Revolution: A History, which inspired Charles Dickens' novel A Tale of Two Cities, argued that European supremacist policies were justified on the grounds that they provided the greatest benefit to "inferior" native peoples. However, even at the time of its publication in 1849, Carlyle's main work on the subject, the Occasional Discourse on the Negro Question, was received poorly by his contemporaries. Before the outbreak of the American Civil War, the Confederate States of America was founded with a constitution that contained clauses which restricted the government's ability to limit or interfere with the institution of "negro" slavery. In the Cornerstone Speech, Confederate vice president Alexander Stephens declared that one of the Confederacy's foundational tenets was white supremacy over black slaves. Following the war, a secret society, the Ku Klux Klan, was formed in the South. Its purpose was to "restore" white supremacy after the Reconstruction period, even though there still was white, Protestant supremacy in the United States, at the time. The group preached supremacy over all other races, as well as supremacy over Jews, Catholics, and other minorities. According to William Nichols, religious antisemitism can be distinguished from modern antisemitism which is based on racial or ethnic grounds. "The dividing line was the possibility of effective conversion ... a Jew ceased to be a Jew upon baptism." However, with racial antisemitism, "Now the assimilated Jew was still a Jew, even after baptism ... . From the Enlightenment onward, it is no longer possible to draw clear lines of distinction between religious and racial forms of hostility towards Jews... Once Jews have been emancipated and secular thinking makes its appearance, without leaving behind the old Christian hostility towards Jews, the new term antisemitism becomes almost unavoidable, even before explicitly racist doctrines appear." One of the first typologies which was used to classify various human races was invented by Georges Vacher de Lapouge (1854–1936), a theoretician of eugenics, who published in 1899 L'Aryen et son rôle social (1899 – "The Aryan and his social role"). In this book, he classifies humanity into various, hierarchized races, spanning from the "Aryan white race, dolichocephalic", to the "brachycephalic", "mediocre and inert" race, best represented by Southern European, Catholic peasants". Between these, Vacher de Lapouge identified the "Homo europaeus" (Teutonic, Protestant, etc.), the "Homo alpinus" (Auvergnat, Turkish, etc.), and finally the "Homo mediterraneus" (Neapolitan, Andalus, etc.) Jews were brachycephalic like the Aryans, according to Lapouge; but exactly for this reason he considered them dangerous; they were the only group, he thought, which was threatening to displace the Aryan aristocracy. Vacher de Lapouge became one of the leading inspirations of Nazi antisemitism and Nazi racist ideology. The Anti-Defamation League and Southern Poverty Law Center condemn writings about "Jewish Supremacism" by Holocaust-denier, former Grand Wizard of the KKK, and conspiracy theorist David Duke as antisemitic – in particular, his book Jewish Supremacism: My Awakening to the Jewish Question. Kevin B. MacDonald, known for his theory of Judaism as a "group evolutionary strategy", has also been accused by the ADL and his own university psychology department of being "antisemitic" and white supremacist in his writings on the subject. Cornel West, an African-American philosopher, writes that Black supremacist religious views arose in America as part of black Muslim theology in response to white supremacism. In Africa, black Southern Sudanese allege that they are subjected to a racist form of Arab supremacy, which they equate with the historic white supremacism of South African apartheid. The alleged genocide and ethnic cleansing in the ongoing War in Darfur has been described as an example of Arab racism. For example, in their analysis of the sources of the conflict, Julie Flint and Alex de Waal say that Colonel Gaddafi, the leader of Libya, sponsored "Arab supremacism" across the Sahara during the 1970s. Gaddafi supported the "Islamic Legion" and the Sudanese opposition "National Front, including the Muslim Brothers and the Ansar, the Umma Party's military wing." Gaddafi tried to use such forces to annex Chad from 1979-81. Gaddafi supported the Sudanese government's war in the South during the early 1980s, and in return he was allowed to use the Darfur region as a "back door to Chad." As a result, the first signs of an "Arab racist political platform" appeared in Darfur in the early 1980s. In Asia, ancient Indians considered all foreigners barbarians. The Muslim scholar Al-Biruni wrote that the Indians called foreigners impure. A few centuries later, Dubois observes that "Hindus look upon Europeans as barbarians totally ignorant of all principles of honour and good breeding... In the eyes of a Hindu, a Pariah (outcaste) and a European are on the same level." The Chinese viewed the Europeans as repulsive, ghost-like creatures, and even as devils. Chinese writers also referred to the Europeans as barbarians. Germany From 1933 to 1945, Nazi Germany, under the rule of Adolf Hitler, promoted the idea of a superior, Aryan Herrenvolk, or master race. The state's propaganda advocated the belief that Germanic peoples, whom they called "Aryans", were a master race or a Herrenvolk that was superior to the Jews, Slavs, and Romani people, so-called "gypsies". Arthur de Gobineau, a French racial theorist and aristocrat, blamed the fall of the ancien régime in France on racial intermixing, which he argued had destroyed the purity of the Nordic race. Gobineau's theories, which attracted a strong following in Germany, emphasized the existence of an irreconcilable polarity between Aryan and Jewish cultures. Fundamentalism Christian Academics Carol Lansing and Edward D. English claim that Christian supremacism was a motivation for the Crusades in the Holy Land, as well as crusades against Muslims and pagans throughout Europe. The blood libel is a widespread European conspiracy theory which led to centuries of pogroms and massacres of European Jewish minorities because it alleged that Jews required the pure blood of a Christian child in order to make matzah for Passover; Thomas of Cantimpré writes of the blood curse which the Jews put upon themselves and all of their generations at the court of Pontius Pilate where Jesus was handed a death sentence: "A very learned Jew, who in our day has been converted to the (Christian) faith, informs us that one enjoying the reputation of a prophet among them, toward the close of his life, made the following prediction: 'Be assured that relief from this secret ailment, to which you are exposed, can only be obtained through Christian blood ("solo sanguine Christiano")." The Atlantic slave trade has also been partially attributed to Christian supremacism. The Ku Klux Klan has been described as a white supremacist Christian organization, as are many other white supremacist groups, such as the Posse Comitatus and the Christian Identity and Positive Christianity movements. Islamic Academics Khaled Abou El Fadl, Ian Lague, and Joshua Cone allege instances of Muslim or Islamic supremacism but they also note that the Qur'an and other Islamic documents always speak of tolerant, protective beliefs, which have been misused, misquoted, and misinterpreted by both Islamic extremists and Islamophobes. Examples of how supremacists have interpreted Islam include the Muslim participation in the African slave trade, the early-20th-century pan-Islamism promoted by Abdul Hamid II, the jizya and rules of marriage in Muslim countries being imposed on non-Muslims, the majority Muslim interpretations of the rules of pluralism in Malaysia, and "defensive" supremacism practiced by some Muslim immigrants in Europe. According to founder of the Muslim Brotherhood, Hassan al-Banna, “it is the nature of Islam to dominate, not to be dominated, to impose its law on all nations and to extend its power to the entire planet”. According to scholar Bernard Lewis, Classical Islamic jurisprudence imposes an open-ended duty on Muslims to expand Muslim rule and Islamic law to all non-Muslims throughout the world. William Kilpatrick in his book The Politically Incorrect Guide to Jihad posits that Islam, unlike other religions, positively commands its adherents to impose its religious law on all peoples, believers and unbelievers alike, whenever possible and by any means necessary. Despite being comparatively more tolerant than Christian Europe there were numerous incidents of massacres and ethnic cleansing of Jews and Christians in North Africa, especially in Morocco, Libya and Algeria where eventually Jews were forced to live in ghettos. Decrees ordering the destruction of synagogues were enacted during the Middle Ages in Egypt, Syria, Iraq and Yemen. At certain times in Yemen, Morocco and Baghdad, Jews were forced to convert to Islam or face death. While there were antisemitic incidents before the 20th century, antisemitism increased dramatically as a result of the Arab–Israeli conflict. After the 1948 Arab–Israeli War, the Palestinian exodus, the creation of the state of Israel and Israeli victories during the wars of 1956 and 1967 were a severe humiliation to Israel's opponents—primarily Egypt, Syria and Iraq. However, by the mid-1970s the vast majority of Jews had left Muslim-majority countries, moving primarily to Israel, France and the United States. The reasons for the exodus are varied and disputed. Jewish Ilan Pappé, an expatriate Israeli historian, writes that the First Aliyah to Israel "established a society based on Jewish supremacy". Joseph Massad, a Professor of Arab Studies, holds that "Jewish supremacism" has always been a "dominating principle" in religious and secular Zionism. Zionism was established with the political goal of creating a sovereign Jewish state where Jews could be the majority, rather than the minority which they were in all nations of the world at that time. Theodor Herzl, the ideological father of Zionism, considered Antisemitism to be an eternal feature of all societies in which Jews lived as minorities, and as a result, he believed that only a separation could allow Jews to escape eternal persecution. "Let them give us sovereignty over a piece of the Earth's surface, just sufficient for the needs of our people, then we will do the rest!" See also Notes Category:Discrimination Category:Ethnic supremacy Category:Fanaticism Category:Narcissism Category:Political theories Category:Prejudices Category:Racism

Wireless mobile terminals utilized in mobile communication systems are increasingly integrated with other electronic devices to expand their functionality and provide a greater range of services to the user. One such electronic device is a geographic position estimator, and in particular, a device that derives position estimates from receiving and decoding signals broadcasted by navigation satellites. An example of a navigation satellite system is GLONASS (Global Navigation Satellite System) developed by the former Soviet Union. This navigation satellite system is similar to the United States Global Positioning System (GPS), the Chinese Compass Navigation System, and the European Union (EU) Galileo Positioning System. Each of these navigation satellite systems broadcasts encoded signals, which may be received at a receiver, such as a geographic position estimator, and may be processed using trilateration, which calculates the receiver's position on the surface of the Earth. To acquire and process a satellite signal requires a receiver to determine data bit boundaries of an incoming stream of data samples. Often a window spanning a single information data bit contains several data samples. Herein, a data sample represents an integration across a PN code period. For example, one data sample represents an integration across a period containing a single 1023-chip PN code. In such cases, the receiving device must determine a starting sample for the window. For example, a window spanning the time of 20 data samples has the possibility of starting at any one of 20 different phases. For meander encoded data samples, a transition rate is effectively doubled by the encoding process, which induces a transition at the center of the data bit (e.g., an information data bit stream has possible transitions at a rate of 50 Hz, whereas a meander encoded data bit stream has possible transitions at a rate of 100 Hz). A meander code is also referred to as a Manchester code. Conventionally for meander encoded information data bits, a two-step process is used to determine a data bit boundary (also referred to as a data bit edge). Prior to this two-step test, the sub-PN code resolution has been determined. That is, the starting point of each repeated PN code has been determined. In a GLONASS system, the PN code repeats every millisecond; therefore, the system has knowledge of the sub-millisecond offset. Thus in GLONASS, the two-step process uses the predetermined 1-ms timing to determine the 10 and 20 ms boundaries. In the first step, initial boundaries at the meander encoded transition rate are determined. The initial boundaries are at a rate twice the rate of the information data bits (i.e., at the meander rate). The initial boundaries both (1) separate neighboring data bits and (2) divide each meander encoded data bit in half. In the second step, the initial boundaries are analyzed to distinguish (1) data bit boundaries separating neighboring data bits from (2) meander encoded center transitions that divide each bit in half. As an example, a window spanning 20 meander encoded data samples at a rate of 1 sample/ms has 20 corresponding possible starting points. First, an initial boundary starting point (or phase) is determined with a 10-ms spacing (100-Hz data rate), then second, a data bit boundary starting point (data bit edge) is determined with a 20-ms spacing (50 Hz data rate). Once the receiving device determines a starting point, the device may remove the meander code and interpret the received data samples. This two-step process has multiple disadvantages and drawbacks. First, when the initial boundary is determined, it may be erroneously set. This guarantees that the second step will also determine an erroneous data bit edge. Additional, the two-step process must complete the first step before beginning the second step. Therefore, an improved method and device for data bit edge detection of meander encoded data samples, for example, from a GLONASS satellite, is desirable.

FIGS. 1 to 3 are drawing illustrating a III-nitride compound semiconductor light emitting device according to the prior art. Specifically, FIG. 1 is a cross-sectional view of the III-nitride compound semiconductor light emitting device, FIG. 2 is an energy band diagram showing that the active layer 13 has an InGaN/GaN multiple-quantum-well structure, and FIG. 3 is an energy band diagram showing that the active layer 13 has an InGaN/InGaN multiple-quantum-well structure. Referring to FIG. 1, to fabricate a nitride compound semiconductor light emitting device according to the prior art, the buffer layer 11, the n-InxAlyGazN layer 12, the active layer 13, the p-InxAlyGazN layer 14 and the transparent electrode layer 15 are successively formed on the substrate 10. Then, mesa etching is performed in such a way to expose the n-InxAlyGazN layer 12. Next, the transparent electrode layers 16 and 17 are formed on the exposed portion of the n-InxAlyGazN layer 12 and the transparent electrode layer 12, respectively. Then, the passivation film 18 is formed on the device. In this respect, x, y and z in InxAlyGazN forming the layers 12 and 14 satisfy the following conditions: x+y+z=1, 0≦x<1, 0≦y<1, and 0<z≦1. As shown in FIG. 2, the active layer 13 has an InGaN/GaN multiple-quantum-well structure having an alternate stacking comprising the InGaN quantum well layers 13b and the GaN barrier layers 13a. Alternatively, as shown in FIG. 3, the active layer 13 has an InGaN/InGaN multiple-quantum-well structure of an alternate stacking comprising the InGaN quantum well layers 13b and the InGaN barrier layers 13a′. In the InGaN/InGaN multiple-quantum-well structure, the barrier layers 13a′ have a lower indium (In) content than that of the quantum well layers 13b. The InGaN/GaN multiple-quantum-well active layer as shown in FIG. 2 has an advantage in that it is more stable in high-current operation or high-temperature operation than InGaN/InGaN. This is because the energy band gap of GaN forming the barrier layers 13a is higher than that of the barrier layers 13a′, leading to an increase in the efficiency of recombination between electrons and holes. However, the InGaN/GaN multiple-quantum-well active layer has a limitation in that the GaN barrier layers 13a should be grown at low temperature since the growth temperature of the InGaN well layers 13b is about 200-350° C. lower than the general growth temperature of high-quality GaN. This makes it difficult to make the GaN barrier layers 13a with high quality levels. Another disadvantage is that the strain of the InGaN well layers 13b and the GaN barrier layers 13a needs to be optimized only with-the thickness of the InGaN well layers 13b and the GaN barrier layers 13a. In the other hand, in the case of the InGaN/InGaN multiple-quantum-well active layer as shown in FIG. 3, the barrier layers 13a′ are also made of InGaN so that the InGaN barrier layers 13a′ with high quality can be made at a growth temperature range similar to that of the InGaN well layers 13b. The InGaN barrier layers 13b made of the same material as the InGaN well layers 13b provide growth surfaces of the next well layers so that the InGaN well layers 13b with high quality can be obtained, leading to a light emitting device having high quantum efficiency. However, the InGaN/InGaN multiple-quantum-well active layer has a disadvantage in that it has low stability in high-current or high-temperature operations as compared to the InGaN/GaN multiple-quantum-well layer. In an attempt to overcome this disadvantage, a p-AlGaN electron-blocking layer is generally formed on the active layer so as to increase the high-temperature stability and efficiency of the active layer. However, this approach has a disadvantage in that, in epitaxial growth of the nitride compound, the deposition of an aluminum (Al)-containing nitride compound, which is difficult to remove easily, occurs within a reactor. This results in a great reduction in mass production. Another disadvantage with the InGaN/InGaN multiple-quantum-well active layer is that the performance and properties of a light emitting diode including the active layer are sensitive to a small change in, the indium composition of the InGaN barrier layer 13a′ thus making the mass production of the diode difficult. Trimethy indium (TMIn) is mainly used as a precursor for indium. It has a melting point of 88° C. and is solid at room temperature. Thus, growing a material with the precursor maintained at a constant composition is relatively difficult compared to other precursors. Particularly, a small change in the indium composition in the barrier layers 13a′ is difficult to detect even by a nondestructive measurement system, such as a PL (photoluminescence) or XRD (X-ray Diffraction) system, and thus, causes a difficulty in quality control and has a great effect on the production of prior products. As described above, the existing typical InGaN/GaN multiple-quantum-well layer and InGaN/InGaN multiple-quantum-well layer have advantages and disadvantages, which are contrary to each other.

Little Laurel Branch Little Laurel Branch, a wildland in the George Washington and Jefferson National Forests of western Virginia, has been recognized by the Wilderness Society as a special place worthy of protection from logging and road construction. The Wilderness Society has designated the area as a “Mountain Treasure”. A strikingly scenic place, Little Laurel Branch is an unusual find in the Cumberland Mountains. It contains part of a water gap, the Breaks of Pine Mountain, with multi-level cliffs and vertical wooded slopes. This wildland is part of the Clinch Ranger District Cluster. Location and access The area is located in the Cumberland Mountains of Southwestern Virginia, about 18 miles southeast of Pikeville, Kentucky and 31 miles west of Richlands, Virginia. It is next to Breaks Interstate Park. There are no trails into the area. However, the Pine Mountain Trail passes within one mile giving an overview into the area at Elkhorn City Overlook, Cave and Skagg Overlook. Skagg Branch Rd. (Rt 2031) passes into the area. The boundary of the wildland as determined by the Wilderness Society is shown in the adjacent map. Additional roads and trails are given on National Geographic Maps 789 (Clinch Ranger District). A great variety of information, including topographic maps, aerial views, satellite data and weather information, is obtained by selecting the link with the wildland's gps coordinates in the upper right of this page. Beyond maintained trails, old logging roads can be used to explore the area. The Cumberland Mountains were extensively timbered in the early twentieth century leaving logging roads that are becoming overgrown but still passable. Old logging roads and railroad grades can be located by consulting the historical topographic maps available from the United States Geological Survey (USGS). The Little Laurel Branch wildarea is covered by USGS topographic map Elkhorn City. Natural history The eastern side of the area is part of the Chimney Cliffs and Russell Fork special biological area, home to a state-listed plant and federally threatened plant. A population of the rare Virgiia spiraea plant is found in the river area. Massive rocks and cliffs profide habitat for calcium rich and iron rich species. Breaks Interstate Park The area is surrounded on three sides by the 4600-acre Breaks Interstate Park. Known as “The Grand Canyon of the South", the park contains a 5 mile long, 25 mile deep gorge, making it a popular visitor destination. Topography The area is part of the Pine and Cumberland Mountains Subsection of the Southern Cumberland Mountain Section of the Central Appalachian Broadleaf Coniferous Forest-Meadow Province. Pine Mountain, a prominent ridge in the region, is a "strike ridge", an escarpment more resistant to weathering than adjacent rock. As time passed, huge blocks of sandstone fell from the ridge leaving a rugged terrain. The Towers, a large rock formation, is found at the eastern end of the area. A section of Russell Fork is considered eligible for the National Wild and Scenic Rivers System. Forest Service management The Forest Service has conducted a survey of their lands to determine the potential for wilderness designation. Wilderness designation provides a high degree of protection from development. The areas that were found suitable are referred to as inventoried roadless areas. Later a Roadless Rule was adopted that limited road construction in these areas. The rule provided some degree of protection by reducing the negative environmental impact of road construction and thus promoting the conservation of roadless areas. Little Laurel Branch was not inventoried in the roadless area review, and therefore not protected from possible road construction and timber sales. The forest service classifies areas under their management by a recreational opportunity setting that informs visitors of the diverse range of opportunities available in the forest. The area includes land designated “Rare Community” and “Geological Area-Unsuitable”. Several gas wells were developed in the area around 1993. There are now possibly four gas leases in the area, some with access roads. See also Clinch Ranger District Cluster References Further reading Stephenson, Steven L., A Natural History of the Central Appalachians, 2013, West Virginia University Press, West Virginia, . Davis, Donald Edward, Where There Are Mountains, An Environmental History of the Southern Appalachians, 2000, University of Georgia Press, Athens, Georgia. . External links George Washington and Jefferson National Forest Wilderness Society Clinch Coalition Category:George Washington and Jefferson National Forests Category:Southwest Virginia

X-rays were discovered in 1895 by Wilhelm Conrad Roentgen (1845-1923) who was a Professor at Wuerzburg University in Germany. Working with a cathode-ray tube in his laboratory, Roentgen observed a fluorescent glow of crystals on a table near his tube. The tube that Roentgen was working with consisted of a glass envelope (bulb) with positive and negative electrodes encapsulated in it. The air in the tube was evacuated, and when a high voltage was applied, the tube produced a fluorescent glow. Roentgen shielded the tube with heavy black paper, and discovered a green colored fluorescent light generated by a material located a few feet away from the tube. He concluded that a new type of ray was being emitted from the tube. This ray was capable of passing through the heavy paper covering and exciting the phosphorescent materials in the room. He found that the new ray could pass through most substances casting shadows of solid objects. Roentgen also discovered that the ray could pass through the tissue of humans, but not bones and metal objects. One of Roentgen's first experiments late in 1895 was a film of the hand of his wife, Bertha. It is interesting that the first use of X-rays were for an industrial (not medical) application, as Roentgen produced a radiograph of a set of weights in a box to show his colleagues. Roentgen's discovery was a scientific bombshell, and was received with extraordinary interest by both scientist and laymen. Scientists everywhere could duplicate his experiment because the cathode tube was very well known during this period. Many scientists dropped other lines of research to pursue the mysterious rays. Newspapers and magazines of the day provided the public with numerous stories, some true, others fanciful, about the properties of the newly discovered rays. Public fancy was caught by this invisible ray with the ability to pass through solid matter, and, in conjunction with a photographic plate, provide a picture of bones and interior body parts. Scientific fancy was captured by the demonstration of a wavelength shorter than light. This generated new possibilities in physics, and for investigating the structure of matter. Much enthusiasm was generated about potential applications of rays as an aid in medicine and surgery. Within a month after the announcement of the discovery, several medical radiographs had been made in Europe and the United States, which were used by surgeons to guide them in their work. In June 1896, only 6 months after Roentgen announced his discovery, X-rays were being used by battlefield physicians to locate bullets in wounded soldiers. Prior to 1912, X-rays were used little outside the realms of medicine and dentistry, though some X-ray pictures of metals were produced. The reason that X-rays were not used in industrial application before this date was because the X-ray tubes (the source of the X-rays) broke down under the voltages required to produce rays of satisfactory penetrating power for industrial purposes. However, that changed in 1913 when the high vacuum X-ray tubes designed by Coolidge became available. The high vacuum tubes were an intense and reliable X-ray source, operating at energies up to 100,000 volts. In 1922, industrial radiography took another step forward with the advent of the 200,000-volt X-ray tube that allowed radiographs of thick steel parts to be produced in a reasonable amount of time. In 1931, General Electric Company developed 1,000,000 volt X-ray generators, providing an effective tool for industrial radiography. That same year, the American Society of Mechanical Engineers (ASME) permitted X-ray approval of fusion welded pressure vessels that further opened the door to industrial acceptance and use. A Second Source of Radiation Shortly after the discovery of X-rays, another form of penetrating rays was discovered. In 1896, French scientist Henri Becquerel discovered natural radioactivity. Many scientists of the period were working with cathode rays, and other scientists were gathering evidence on the theory that the atom could be subdivided. Some of the new research showed that certain types of atoms disintegrate by themselves. It was Henri Becquerel who discovered this phenomenon while investigating the properties of fluorescent minerals. Becquerel was researching the principles of fluorescence, wherein certain minerals glow (fluoresce) when exposed to sunlight. He utilized photographic plates to record this fluorescence. One of the minerals Becquerel worked with was a uranium compound. On a day when it was too cloudy to expose his samples to direct sunlight, Becquerel stored some of the compound in a drawer with his photographic plates. Later when he developed these plates, he discovered that they were fogged (exhibited exposure to light). Becquerel questioned what would have caused this fogging. He knew he had wrapped the plates tightly before using them, so the fogging was not due to stray light. In addition, he noticed that only the plates that were in the drawer with the uranium compound were fogged. Becquerel concluded that the uranium compound gave off a type of radiation that could penetrate heavy paper and expose photographic film. Becquerel continued to test samples of uranium compounds and determined that the source of radiation was the element uranium. Bacquerel's discovery was, unlike that of the X-rays, virtually unnoticed by laymen and scientists alike. Relatively few scientists were interested in Becquerel's findings. It was not until the discovery of radium by the Curies two years later that interest in radioactivity became widespread. While working in France at the time of Becquerel's discovery, Polish scientist Marie Curie became very interested in his work. She suspected that a uranium ore known as pitchblende contained other radioactive elements. Marie and her husband, French scientist Pierre Curie, started looking for these other elements. In 1898, the Curies discovered another radioactive element in pitchblende, and named it 'polonium' in honor of Marie Curie's native homeland. Later that year, the Curies discovered another radioactive element which they named radium, or shining element. Both polonium and radium were more radioactive than uranium. Since these discoveries, many other radioactive elements have been discovered or produced. Radium became the initial industrial gamma ray source. The material allowed castings up to 10 to 12 inches thick to be radiographed. During World War II, industrial radiography grew tremendously as part of the Navy's shipbuilding program. In 1946, man-made gamma ray sources such as cobalt and iridium became available. These new sources were far stronger than radium and were much less expensive. The manmade sources rapidly replaced radium, and use of gamma rays grew quickly in industrial radiography. Health Concerns The science of radiation protection, or "health physics" as it is more properly called, grew out of the parallel discoveries of X-rays and radioactivity in the closing years of the 19th century. Experimenters, physicians, laymen, and physicists alike set up X-ray generating apparatuses and proceeded about their labors with a lack of concern regarding potential dangers. Such a lack of concern is quite understandable, for there was nothing in previous experience to suggest that X-rays would in any way be hazardous. Indeed, the opposite was the case, for who would suspect that a ray similar to light but unseen, unfelt, or otherwise undetectable by the senses would be damaging to a person? More likely, or so it seemed to some, X-rays could be beneficial for the body. Inevitably, the widespread and unrestrained use of X-rays led to serious injuries. Often injuries were not attributed to X-ray exposure, in part because of the slow onset of symptoms, and because there was simply no reason to suspect X-rays as the cause. Some early experimenters did tie X-ray exposure and skin burns together. The first warning of possible adverse effects of X-rays came from Thomas Edison, William J. Morton, and Nikola Tesla who each reported eye irritations from experimentation with X-rays and fluorescent substances. Today, it can be said that radiation ranks among the most thoroughly investigated causes of disease. Although much still remains to be learned, more is known about the mechanisms of radiation damage on the molecular, cellular, and organ system than is known for most other health stressing agents. Indeed, it is precisely this vast accumulation of quantitative dose-response data that enables health physicists to specify radiation levels so that medical, scientific, and industrial uses of radiation may continue at levels of risk no greater than, and frequently less than, the levels of risk associated with any other technology. X-rays and Gamma rays are electromagnetic radiation of exactly the same nature as light, but of much shorter wavelength. Wavelength of visible light is on the order of 6000 angstroms while the wavelength of x-rays is in the range of one angstrom and that of gamma rays is 0.0001 angstrom. This very short wavelength is what gives x-rays and gamma rays their power to penetrate materials that light cannot. These electromagnetic waves are of a high energy level and can break chemical bonds in materials they penetrate. If the irradiated matter is living tissue, the breaking of chemical bonds may result in altered structure or a change in the function of cells. Early exposures to radiation resulted in the loss of limbs and even lives. Men and women researchers collected and documented information on the interaction of radiation and the human body. This early information helped science understand how electromagnetic radiation interacts with living tissue. Unfortunately, much of this information was collected at great personal expense.

As the water-intensive practice of fracking continues to spread, the amount of wastewater being buried in disposal wells around Texas has skyrocketed. But the wells bring concerns about truck traffic and the possibility of groundwater contamination. In drilling regions like the Permian Basin, where the water needs of fracking have run up against a historic drought, drillers are increasingly turning to brackish groundwater previously thought too expensive to use. Most fracking operations use several million gallons of water. But with water increasingly scarce and costly around Texas, a few companies have begun using alternative liquids, such as propane. Experts say the technology still has far to go. In Pennsylvania's Marcellus Shale, natural gas companies recycle water as a matter of course. But recycling is only getting started in the Texas oilfields because using freshwater for hydraulic fracturing is cheap. The amount of water used in hydraulic fracturing has stirred concerns around Texas, especially as the drought wears on. Aware that they are under the spotlight, drillers are testing out recycling and other water-saving techniques.