metadata
library_name: diffusers
pipeline_tag: text-to-image
inference: true
base_model: stabilityai/sd-turbo
DPO LoRA Stable Diffusion v2-1
Model trained with LoRA implementation of Diffusion DPO Read more here
Base Model: https://huggingface.co/stabilityai/stable-diffusion-2-1
Running
from diffusers import DiffusionPipeline
from diffusers.utils import make_image_grid
import torch
pipe = DiffusionPipeline.from_pretrained(
"stabilityai/sd-turbo",
# "stabilityai/stable-diffusion-2-1",
torch_dtype=torch.float16, variant="fp16"
)
pipe.load_lora_weights("radames/dpo-lora-sd2.1")
seed = 123123
pipe.to("cuda")
prompt = "portrait headshot professional of elon musk"
negative_prompt = "3d render, cartoon, drawing, art, low light"
generator = torch.Generator().manual_seed(seed)
images = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
width=512,
height=512,
num_inference_steps=2,
generator=generator,
guidance_scale=1.0,
num_images_per_prompt=4
).images
make_image_grid(images, 1, 4)