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A newer version of the Gradio SDK is available:
5.9.1
ControlNet training example for Stable Diffusion XL (SDXL)
The train_controlnet_sdxl.py
script shows how to implement the ControlNet training procedure and adapt it for Stable Diffusion XL.
Running locally with PyTorch
Installing the dependencies
Before running the scripts, make sure to install the library's training dependencies:
Important
To make sure you can successfully run the latest versions of the example scripts, we highly recommend installing from source and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install -e .
Then cd in the examples/controlnet
folder and run
pip install -r requirements_sdxl.txt
And initialize an 🤗Accelerate environment with:
accelerate config
Or for a default accelerate configuration without answering questions about your environment
accelerate config default
Or if your environment doesn't support an interactive shell (e.g., a notebook)
from accelerate.utils import write_basic_config
write_basic_config()
When running accelerate config
, if we specify torch compile mode to True there can be dramatic speedups.
Circle filling dataset
The original dataset is hosted in the ControlNet repo. We re-uploaded it to be compatible with datasets
here. Note that datasets
handles dataloading within the training script.
Training
Our training examples use two test conditioning images. They can be downloaded by running
wget https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/controlnet_training/conditioning_image_1.png
wget https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/controlnet_training/conditioning_image_2.png
Then run huggingface-cli login
to log into your Hugging Face account. This is needed to be able to push the trained ControlNet parameters to Hugging Face Hub.
export MODEL_DIR="stabilityai/stable-diffusion-xl-base-1.0"
export OUTPUT_DIR="path to save model"
accelerate launch train_controlnet_sdxl.py \
--pretrained_model_name_or_path=$MODEL_DIR \
--output_dir=$OUTPUT_DIR \
--dataset_name=fusing/fill50k \
--mixed_precision="fp16" \
--resolution=1024 \
--learning_rate=1e-5 \
--max_train_steps=15000 \
--validation_image "./conditioning_image_1.png" "./conditioning_image_2.png" \
--validation_prompt "red circle with blue background" "cyan circle with brown floral background" \
--validation_steps=100 \
--train_batch_size=1 \
--gradient_accumulation_steps=4 \
--report_to="wandb" \
--seed=42 \
--push_to_hub
To better track our training experiments, we're using the following flags in the command above:
report_to="wandb
will ensure the training runs are tracked on Weights and Biases. To use it, be sure to installwandb
withpip install wandb
.validation_image
,validation_prompt
, andvalidation_steps
to allow the script to do a few validation inference runs. This allows us to qualitatively check if the training is progressing as expected.
Our experiments were conducted on a single 40GB A100 GPU.
Inference
Once training is done, we can perform inference like so:
from diffusers import StableDiffusionXLControlNetPipeline, ControlNetModel, UniPCMultistepScheduler
from diffusers.utils import load_image
import torch
base_model_path = "stabilityai/stable-diffusion-xl-base-1.0"
controlnet_path = "path to controlnet"
controlnet = ControlNetModel.from_pretrained(controlnet_path, torch_dtype=torch.float16)
pipe = StableDiffusionXLControlNetPipeline.from_pretrained(
base_model_path, controlnet=controlnet, torch_dtype=torch.float16
)
# speed up diffusion process with faster scheduler and memory optimization
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
# remove following line if xformers is not installed or when using Torch 2.0.
pipe.enable_xformers_memory_efficient_attention()
# memory optimization.
pipe.enable_model_cpu_offload()
control_image = load_image("./conditioning_image_1.png").resize((1024, 1024))
prompt = "pale golden rod circle with old lace background"
# generate image
generator = torch.manual_seed(0)
image = pipe(
prompt, num_inference_steps=20, generator=generator, image=control_image
).images[0]
image.save("./output.png")
Notes
Specifying a better VAE
SDXL's VAE is known to suffer from numerical instability issues. This is why we also expose a CLI argument namely --pretrained_vae_model_name_or_path
that lets you specify the location of an alternative VAE (such as madebyollin/sdxl-vae-fp16-fix
).
If you're using this VAE during training, you need to ensure you're using it during inference too. You do so by:
+ vae = AutoencoderKL.from_pretrained(vae_path_or_repo_id, torch_dtype=torch.float16)
controlnet = ControlNetModel.from_pretrained(controlnet_path, torch_dtype=torch.float16)
pipe = StableDiffusionXLControlNetPipeline.from_pretrained(
base_model_path, controlnet=controlnet, torch_dtype=torch.float16,
+ vae=vae,
)